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1027433 | https://en.wikipedia.org/wiki/Wiwaxia | Wiwaxia | Wiwaxia is a genus of soft-bodied animals that were covered in carbonaceous scales and spines that protected it from predators. Wiwaxia fossils—mainly isolated scales, but sometimes complete, articulated fossils—are known from early Cambrian and middle Cambrian fossil deposits across the globe. The living animal would have measured up to when fully grown, although a range of juvenile specimens are known, the smallest being long.
Wiwaxia affinity has been a matter of debate: Researchers were long split between two possibilities. On the one hand, its rows of scales looked superficially similar to certain scale worms (annelids); conversely, its mouthparts and general morphology suggested a relationship to the shell-less molluscs. More recently, evidence for a molluscan affinity has been accumulating, based on new details of Wiwaxia mouthparts, scales, and growth history. The name derives from Wiwaxy Peak in British Columbia.
The proposed clade Halwaxiida contains Wiwaxia as well as several similar Cambrian animals.
Description
This article concentrates on the species Wiwaxia corrugata, which is known from hundreds of complete specimens in the Burgess Shale; other species are known only from fragmentary material or limited sample sizes.
Body
Wiwaxia was bilaterally symmetrical; viewed from the top the body was elliptical with no distinct head or tail, and from the front or rear it was almost rectangular. It reached in length. Estimating their height is difficult because specimens were compressed after death; a typical specimen may have been high excluding the spines on their backs. The ratio of width to length does not appear to change as the animals grew.
Wiwaxia flat underside was soft and unarmored; most of the surface was occupied by a single slug-like foot. Little is known of the internal anatomy, although the gut apparently ran straight and all the way from the front to the rear. At the front end of the gut, about from the animal's front in an average specimen about long, there was a feeding apparatus that consisted of two (or in rare large specimens three) rows of backward-pointing conical teeth. The feeding apparatus was tough enough to be frequently preserved, but unmineralized and fairly flexible.
Sclerites
The animal was covered in eight rows of small ribbed armor plates called sclerites; these lay flat against the body, overlapped so that the rear of one covered the front of the one behind, and formed five main regions—the top; the upper part of the sides; the lower part of the sides; the front; and the bottom. Most of the sclerites were shaped like oval leaves, but the ventro-lateral ones, nearest the sea-floor, were crescent-shaped, rather like flattened bananas, and formed a single row. Larger specimens (>~15 mm) bear two rows of ribbed spines running from front to rear, one along each side of the top surface, and projecting out and slightly upward, with a slight upward curve near the tips. Although the spines in the middle of each row are usually the longest, up to long, a few specimens have rather short middle spines that represent part-grown replacements.
Each sclerite was rooted separately in the body; the roots of body sclerites are 40% of the external length or a little less, while the roots of the spines are a little over 25% of the external length; all were rooted in pockets in the skin, rather like the follicles of mammalian hair. The roots of the body sclerites were significantly narrower than the sclerites, but the spines had roots about as wide as their bases; both types of root were made of fairly soft tissue. They bore protrusive, presumably structural, ribs on their upper and (seemingly) lower surfaces. The sclerites and spines were not mineralized, but made of a tough organic (carbon-based) biopolymer. Butterfield (1990) examined some sclerites under both optical and scanning electron microscopes and concluded that they were not hollow, and that the bases split and spread to form the blades, a pattern that is also seen in monocot leaves. The sclerites bear an internal fabric of longitudinal chambers, which suggest that they were secreted from their bases in the manner of Lophotrochozoan sclerites.
As known from Marrella and Canadia, its sclerites may have been iridescent due to evidence of diffraction grating, although questioned in later study.
Ontogeny
Wiwaxia scleritome comprises eight rows of sclerites, arranged in bundles. Sclerites are periodically shed and replaced during growth, with the number of sclerites in a given bundle increasing as the animal ages to produce a thicker scleritome. Once specimens reach a certain size, spines are added to the scleritome; this size is ~15 mm in W. corrugata but substantially smaller in W. taijengensis. One juvenile specimen was originally interpreted as molting, but in fact represents a single, folded, individual.
Ecology
The long dorsal spines may have been a defense against predators. Wiwaxia apparently moved by contractions of a slug-like foot on its underside. In one specimen a small brachiopod, Diraphora bellicostata, appears to be attached to one of the ventro-lateral sclerites. This suggests that adult Wiwaxia did not burrow or even plough much into the sea-floor as they moved. Two other specimens of Diraphora bellicostata have been found attached to dorsal sclerites. Wiwaxia appears to have been solitary rather than gregarious. The feeding apparatus may have acted as a rasp to scrape bacteria off the top of the microbial mat that covered the sea-floor, or as a rake to gather food particles from the sea-floor.
Classification
During the Cambrian, most of the main groupings of animals recognised today were beginning to diverge. Consequently, many lineages (that would later become extinct) appear intermediate to two or more modern groups, or lack features common to all modern members of a group, and hence fall into the "stem group" of a modern taxon. Debate is ongoing as to whether Wiwaxia can be placed within a modern crown group and, if it cannot, in which group's stem it falls. When Walcott first described Wiwaxia, he regarded it as a polychaete annelid worm, and its sclerites as similar to the elytra ("scales") of annelids. More recently the debate has been intense, and proposed classifications include: a member of an extinct phylum distantly related to the molluscs; a crown-group polychaete; a stem-group annelid; a problematic bilaterian; a stem- or possibly primitive crown-group mollusc.
In 1985, Simon Conway Morris agreed that there were similarities to polychaetes, but considered that Wiwaxia sclerites were different in construction from annelids' elytra. He was more impressed by the similarities between Wiwaxia feeding apparatus and a molluscan radula, and assigned the animal to a new taxon Molluscata, which he proposed should also contain the molluscs and hyolithids. When he later described the first fairly complete specimens of Halkieria, he suggested that these were closely related to Wiwaxia.
Nick Butterfield, then a postgraduate paleontologist at Harvard inspired by Stephen Jay Gould's lectures, agreed that the sclerites were not like elytra, which are relatively fleshy and soft. However, since the sclerites were solid, he concluded that Wiwaxia could not be a member of the "Coeloscleritophora", a taxon that had been proposed in order to unite organisms with hollow sclerites, and could not be closely related to the halkieriids, which have hollow sclerites. Instead he thought that they were very similar in several ways to the chitinous bristles (setae) that project from the bodies of modern annelids and in some genera form leaf-like scales that cover the back like roof tiles—in composition, in detailed structure, in how they were attached to the body via "follicles" and in overall appearance. Some modern annelids also develop on each side rows of longer bristles, which both Walcott and Butterfield considered similar to Wiwaxia dorsal spines. including the halkieriids.
Butterfield also contended that Wiwaxia feeding apparatus, instead of being mounted in the middle of its "head", was just as likely to be mounted in two parts on the sides of the "head", an arrangement that is common in polychaetes. He went so far as to classify Wiwaxia as a member of a modern order, Phyllodocida, and pointed out that Wiwaxia lack of obvious segmentation is no barrier to this, as some modern polychaetes also show no segmentation except during development. He later noted that Wiwaxia lack some polychaete features which he would expect to be easily preserved in fossils, and therefore a stem-group annelid, in other words an evolutionary "aunt" of modern annelids.
Conway Morris and Peel (1995) largely accepted Butterfield's arguments and treated Wiwaxia as an ancestor or "aunt" of the polychaetes, and said Butterfield had informed them that the microscopic structure of Wiwaxias sclerites was identical to that of the bristles of two Burgess Shale polychaetes Burgessochaeta and Canadia. Conway Morris and Peel also wrote that one specimen of Wiwaxia showed traces of a small shell, possibly a vestige left over from an earlier stage in the animal's evolution, and noted that one group of modern polychaetes also has what may be a vestigial shell. However, they maintained that Wiwaxia feeding apparatus was much more like a molluscan radula. They also argued that Wiwaxia was fairly closely related to and in fact descended from the halkieriids, as the sclerites are divided into similar groups, although those of halkieriids were much smaller and more numerous; they also said that in 1994 Butterfield had found Wiwaxia sclerites that were clearly hollow. They presented a large cladogram according to which:
The earliest halkieriids were a "sister" group to the molluscs, in other words descendants of a fairly closely related common ancestor.
The halkieriids which Conway Morris had found in Greenland's Sirius Passet lagerstätte were a "sister" group to brachiopods, animals whose modern forms have bivalve shells, but differ from molluscs in having muscular stalks and a distinctive feeding apparatus, the lophophore.
Another halkieriid genus, Thambetolepis, was a "great aunt" of annelids and Wiwaxia was an "aunt" of annelids.
Marine biologist Amélie H. Scheltema et al. (2003) argued that Wiwaxia feeding apparatus is very similar to the radulas of some modern shell-less aplacophoran molluscs, and that the sclerites of the two groups are very similar. They concluded that Wiwaxia was a member of a clade that includes molluscs. Scheltema has also highlighted similarities between Wiwaxia and the larvae of certain solenogaster molluscs, which bear iterated calcareous sclerites arranged into three symmetrical lateral zones.
Danish zoologist Danny Eibye-Jacobsen argued in 2004 that Wiwaxia lacks any characters that would firmly place it as a polychaete or annelid. Eibye-Jacobsen regarded bristles as a feature shared by molluscs, annelids and brachiopods. Hence even if Wiwaxia sclerites closely resembled bristles, which he doubted, this would not prove that Wiwaxia closest relative was the annelids. He also pointed out that the very different numbers of sclerites in the various zones of Wiwaxia body do not correspond to any reasonable pattern of segmentation; while Eibye-Jacobsen did not think that this alone would prevent classification of Wiwaxia as a polychaete, he thought it was a serious objection given the lack of other clearly polychaete features. In his opinion there were no strong grounds for classifying Wiwaxia as a proto-annelid or a proto-mollusc, although he thought the objections against classification as a proto-annelid were the stronger.
Butterfield returned to the debate in 2006, repeating the arguments he presented in 1990 for regarding Wiwaxia as an early polychaete and adding that, while bristles are a feature of several groups, they appear as a covering over the back only in polychaetes.
A 2012 study redescribing the mouthparts found a number of similarities with the molluscan radula, and overthrew some of the better arguments for an annelid affinity, seemingly demonstrating that Wiwaxia was indeed a mollusc.
Occurrence
Wiwaxia was originally described by G. F. Matthew in 1899, from an isolated spine that had been found earlier in the Ogyopsis Shale, and classified as "Orthotheca corrugata". Further specimens were found by American paleontologist Charles Doolittle Walcott in 1911 as a result of one of his field trips to the nearby Burgess Shale (Miaolingian, ) in the Canadian Rocky Mountains; he classified it as a member of the polychaete group of annelid worms in its own genus Wiwaxia corrugata, citing similarities to the Aphroditidae and Polynoidae.
In 1966 and 1967, a team led by Harry B. Whittington revisited the Burgess Shale and found so many fossils that it took years to analyze them all. Four hundred and sixty-four complete specimens of Wiwaxia are known from the Greater Phyllopod bed, where they comprise 0.88% of the community. Eventually in 1985 Simon Conway Morris, then a member of Whittington's team, published a detailed description and concluded that Wiwaxia was not a polychaete. All the known specimens came from in and around the Burgess Shale until 1991, when fragmentary fossils were reported from Australia's Georgina Basin. In 2004 additional finds which may represent two different species were reported from the same area.
Articulated specimens are known from Cambrian Stage 3 of Xiaoshiba, China; fragmentary specimens have also been found in the Chengjiang, Cambrian Series 2 deposits in Guizhou, China, and in China's lowermost Miaolingian beds of the Kaili Formation, in the Middle Cambrian beds of the Tyrovice Member, Buchava Formation of the Czech Republic, in the Lower Cambrian Mount Cap formation (Mackenzie Mountains, Canada), in the Emu Bay Shale of Kangaroo Island, South Australia, upper Botomian Stage of the Lower Cambrian, and in the Middle Botomian Sinsk Biota of Siberia, Russia. Isolated sclerites are also common in the small carbonaceous fossil record. Taken together, these finds show that Wiwaxia had a truly cosmopolitan distribution, occurring at all palaeolatitudes and on most palaeocontinents. The Chinese material was originally considered to represent a separate species; like W. corrugata, it possessed spines and regions of sclerites (although it is only known from disarticulated remains), but the sclerites bear a higher density of ribs, and there are two distinct thicknesses of rib (i.e. larger and smaller). At a microscopic level, the sclerites do not differ from Burgess Shale or Mount Cap sclerites, but the Chinese material seems to have developed spines from an early age, distinguishing it from the W. corrugata. The knob-bearing sclerites from all three localities seem to belong to a different species, and a further species is represented in the Xiaoshiba deposits. What is surprising is the limited variety exhibited between species: all have a fundamentally equivalent scleritome, displaying a notable degree of morphological stasis for some 15 Ma.
Isolated spines are more common than sclerites in localities with a poor preservation potential, suggesting that the spines were more recalcitrant (or more commonly collected); however, in well-preserved sites such as the Phyllopod bed, spine and sclerite abundance is comparable, in disarticulated instances, to the proportions on complete fossils.
Younger spines of possible wiwaxiid origin have been observed from the Valongo Formation (Middle Ordovician: Dapingian-Darriwilian) of northern Portugal and have been reported, if not described, from the Fezouata Biota.
| Biology and health sciences | Mollusks | Animals |
1028100 | https://en.wikipedia.org/wiki/Sit-up | Sit-up | The sit-up is an abdominal endurance training exercise to strengthen, tighten and tone the abdominal muscles. It is similar to a curl-up (that target the rectus abdominis and also work the external and internal obliques), but sit-ups have a fuller range of motion and condition additional muscles.
Form
Sit-ups begin with the practicing individual lying with their back on the floor. Typically, this is done with the arms across the chest or hands behind the head. The knees and toes are bent to reduce stress on the back muscles and spine. Both the upper and lower vertebrae are elevated from the floor until everything superior to the buttocks is not touching the ground. Some argue that sit-ups can be dangerous due to high compressive lumbar load and may be replaced with the crunch in exercise programs. Performing alternative abdominal exercises to sit-ups actually increases the ability to do sit-ups.
Performing sit-ups do not cause the spot reduction of fat at the waist.
Gaining a "six pack" requires both abdominal muscle hypertrophy training and fat loss over the abdomen—which can only be done by losing fat from the body as a whole.
Variations
The movement can be made easier by placing the arms further down away from the head. Typical variations to this include crossing the arms to place the palms on the front of the shoulders and extending the arms down to the sides with palms on the floor. The 'arms on shoulders' variation is also used to make the incline sit-up easier.
More intense movement is achieved by doing weighted sit-ups, incline sit-ups with arms behind neck and even harder by doing the weighted incline sit-up.
Health risks
With improper form, full sit-ups have been found to cause back pain and arching of the lower back, increasing the risk of back injury.
In 2015, it was revealed that every branch of the U.S. armed forces have begun to phase out sit-ups and crunches, due to the high rates of lower-back injury. They have been replaced by planks.
| Biology and health sciences | Physical fitness | Health |
1028265 | https://en.wikipedia.org/wiki/Asterism%20%28astronomy%29 | Asterism (astronomy) | An asterism is an observed pattern or group of stars in the sky. Asterisms can be any identified pattern or group of stars, and therefore are a more general concept than the 88 formally defined constellations. Constellations are based on asterisms, but unlike asterisms, constellations outline and today completely divide the sky and all its celestial objects into regions around their central asterisms. For example, the asterism known as the Big Dipper or the Plough comprises the seven brightest stars in the constellation Ursa Major. Another asterism is the triangle, within the constellation of Capricornus.
Asterisms range from simple shapes of just a few stars to more complex collections of many stars covering large portions of the sky. The stars themselves may be bright naked-eye objects or fainter, even telescopic, but they are generally all of a similar brightness to each other. The larger brighter asterisms are useful for people who are familiarizing themselves with the night sky.
The patterns of stars seen in asterisms are not necessarily a product of any physical association between the stars, but are rather the result of the particular perspectives of their observations. For example the Summer Triangle is a purely observational physically unrelated group of stars, but the stars of Orion's Belt are all members of the Orion OB1 association and five of the seven stars of the Big Dipper are members of the Ursa Major Moving Group. Physical associations, such as the Hyades or Pleiades, can be asterisms in their own right and part of other asterisms at the same time.
Background of asterisms and constellations
In many early civilizations, it was common to associate groups of stars in connect-the-dots stick-figure patterns. Some of the earliest records are those of ancient India in the Vedanga Jyotisha and the Babylonians. Different cultures identified different constellations, although a few of the more obvious patterns tend to appear in the constellations of multiple cultures, such as those of Orion and Scorpius. As anyone could arrange and name a grouping of stars there was no distinct difference between a constellation and an asterism. For example, Pliny the Elder mentions 72 asterisms in his book Naturalis Historia.
A general list containing 48 constellations likely began to develop with the astronomer Hipparchus (c. 190 – c. 120 BCE). As constellations were considered to be composed only of the stars that constituted the figure, it was always possible to use any leftover stars to create and squeeze in a new grouping among the established constellations.
Exploration by Europeans to other parts of the globe exposed them to stars previously unknown to them. Two astronomers particularly known for greatly expanding the number of southern constellations were Johann Bayer (1572–1625) and Nicolas Louis de Lacaille (1713–1762). Bayer had listed twelve figures made out of stars that were too far south for Ptolemy to have seen. Lacaille created 14 new groups, mostly for the area surrounding South Celestial Pole. Many of these proposed constellations have been formally accepted, but the rest have remained as asterisms.
In 1928, the International Astronomical Union (IAU) precisely divided the sky into 88 official constellations following geometric boundaries encompassing all of the stars within them. Any additional new selected groupings of stars or former constellations are often considered as asterisms. However, technical distinctions between the terms 'constellation' and 'asterism' often remain somewhat ambiguous.
Asterisms consisting of first-magnitude stars
Some asterisms consist completely of bright first-magnitude stars, which mark out simple geometric shapes.
The Summer Triangle of Deneb, Altair, and Vega – α Cygni, α Aquilae, and α Lyrae – is prominent in the northern hemisphere summer skies, as its three stars are all of the 1st magnitude. The stars of the Triangle are in the band of the Milky Way which marks the galactic equator, and are in the direction of the Galactic Center.
The Winter Triangle is visible in the northern sky's winter and comprises the first magnitude stars Betelgeuse, Sirius and Procyon (the second and fourth closest star or star system visible without aid).
The larger northern Winter Hexagon includes seven of the twenty-two first-magnitude stars visible in the sky, with Pollux, Capella, Aldebaran, Rigel, Sirius and Procyon, and with the 2nd-magnitude Castor on the periphery, and Betelgeuse off-center. Adding Betelgeuse then it is known as the Heavenly 'G'''. It encircles the galactic anticenter, as well as incorporates constellations such as Gemini and Orion. It also includes in the background of Aldebaran the Hyades, the nearest star cluster and one of five first-magnitude deep-sky objects, two of which can be seen just north-east of the Hyades, the Pleiades also in the Taurus constellation and the Alpha Persei Cluster (with Alcyone and Mirfak as the brightest stars).
The northern Spring Triangle consists of Arcturus, Regulus and Spica.
The Great Diamond consisting of Arcturus, Spica, Denebola and Cor Caroli, the latter two not being first-magnitude stars. An east-west line from Arcturus to Denebola forms an equilateral triangle with Cor Caroli to the North, and another with Spica to the South. Together these two triangles form the Diamond. Formally, the stars of the Diamond are in the constellations Boötes, Virgo, Leo, and Canes Venatici.
Other asterisms consist partially of multiple first-magnitude stars.
The Southern Cross including the first-magnitude stars Acrux and Mimosa, west of the Carina Nebula (one of five first-magnitude deep-sky objects), and with the first-magnitude stars Alpha Centauri (the closest star to the Sun) and Beta Centauri pointing at the cross, distinguishing the cross from less bright and similar asterisms like the Diamond Cross or False Cross.
All other first-magnitude stars are the only such stars in their asterisms or constellations, with Canopus in the Argo Navis asterism south of Sirius, visually east of the Carina Nebula and near the Large Magellanic Cloud (both being first-magnitude deep-sky objects), Achernar in the Eridanus constellation east of Canopus, Fomalhaut in the Southern Fish constellation east of Achernar and Antares in the Scorpius constellation visually near the Galactic Center.
Constellation-based asterisms
The Big Dipper, also known as The Plough or Charles's Wain, is composed of the seven brightest stars in Ursa Major. These stars delineate the Bear's hindquarters and exaggerated tail, or alternatively, the "handle" forming the upper outline of the bear's head and neck. With its longer tail, Ursa Minor hardly appears bearlike at all, and is widely known by its pseudonym, the Little Dipper.
The Northern Cross in Cygnus. The upright runs from Deneb (α Cyg) in the Swan's tail to Albireo (β Cyg) in the beak. The transverse runs from ε Cygni in one wing to δ Cygni in the other.
The Southern Cross is an asterism by name, but the whole area is now recognised as the constellation Crux. The main stars are Alpha, Beta, Gamma, Delta, and arguably also Epsilon Crucis. Earlier, Crux was deemed an asterism when Bayer created it in Uranometria (1603) from the stars in the hind legs of Centaurus, decreasing the size of Centaur. These same stars were probably identified by Pliny the Elder in his Naturalis Historia as the asterism 'Thronos Caesaris.'
The Fish Hook is the traditional Hawaiian name for Scorpius. The image will be even more obvious if the chart's lines from Antares (α Sco) to Beta Scorpii (β Sco) and Pi Scorpii (π Sco) are replaced with a line from Beta through Delta Scorpii (δ Sco) to Pi forming a large capped "J." Adding vertical lines to connect the limbs at the left and right in the main diagram of Hercules will complete the figure of the Butterfly.
Boötes is sometimes known as the Ice Cream Cone. It is also known as the Kite.
The stars of Cassiopeia form a W which is often used as a nickname.
The Great Square of Pegasus is the quadrilateral formed by the stars Markab, Scheat, Algenib, and Alpheratz, representing the body of the winged horse. The asterism was recognized as the constellation ASH.IKU "The Field" on the MUL.APIN cuneiform tablets from about 1100 to 700 BC. Alpheratz is now only considered a part of the constellation Andromeda whereas formerly the star was a part of both constellations.
The Bowl of Virgo is formed by the stars Beta, Gamma, Delta, Epsilon and Eta Virginis. Together with Spica, they form a Y shape.
The Three Leaps of the Gazelle consists of three pairs of stars in Ursa Major aligned in a row spanning about 30 degrees. In Arabic lore, the star pairs are pictured as the hoof prints of a gazelle startled from a pond by Leo the lion. (The "pond" is pictured as the Coma Star Cluster.) The first pair of stars are Xi and Nu, second pair Upsilon and Lambda, third pair Kappa and Iota Ursa Majoris. The pairs also mark three of the bear's paws.
Some asterisms refer to portions of traditional constellation figures. These include:
The Water Jar or Urn of Aquarius is a Y-shaped figure centered upon Zeta Aquarii and includes Gamma, Eta and Pi. It pours water in a stream of more than 20 stars terminating with the star Fomalhaut.
The Crab Breast of Cancer is a quadrilateral formed by the four stars Gamma, Delta, Eta and Theta Cancri which make up the carapace (inner shell) of the Crab. Contained within is the Beehive Cluster (Messier 44) which includes Epsilon Cancri.
The Snake Head is the westernmost portion of Hydra consisting of the stars Delta, Epsilon, Zeta, Eta, Rho and Sigma Hydrae.
Orion's Belt consists of the three bright stars Zeta (Alnitak), Epsilon (Alnilam) and Delta Orionis (Mintaka) which form the belt of Orion.
The Bull's Face of Taurus is a V-shaped figure formed by prominent members of the Hyades cluster, including stars Gamma, Delta¹, Delta², Delta³, Epsilon, Theta Tauri, as well as the bright star Alpha Tauri (Aldebaran) which forms the red eye of the Bull.
Other particular asterisms
Other asterisms are also composed of stars from one constellation, but do not refer to the traditional figures.
Four stars (Beta, Upsilon, Theta, and Omega Carinae) form a well-shaped diamond – the Diamond Cross.
The Saucepan or Pot, being the same stars as the Belt and Sword of Orion. The end of the handle is at ι Orionis, with the far rim at η Orionis.
The four central stars in Hercules, Epsilon (ε Her), Zeta (ζ Her), Eta (η Her), and Pi (π Her), form the Keystone. The bright globular cluster Messier 13 lies along the western segment, between Zeta and Eta.
The curve of stars at the front end of the Lion from Epsilon (ε Leo) to Regulus (α Leo), looking much like a mirror-image question mark, has long been known as the Sickle.
The brighter stars of Sagittarius form the Teapot. (The Large Sagittarius Star Cloud appears to be steam emerging from the "spout".)
Northeast of the Teapot asterism lies the fainter Teaspoon, consisting of the stars ξ¹, ξ², ο, π, ρ¹ and ρ² Sagitarii.
Four bright stars in Delphinus (Sualocin or α Delphini, Rotanev or β Delphini, γ Delphini and δ Delphini) form Job's Coffin.
The Terebellum is a small quadrilateral of four faint stars (Omega, 59, 60, 62) in Sagittarius' hindquarters.
Just south of Pegasus, the western fish of Pisces is home to the Circlet formed from Gamma (γ Piscium), Kappa (κ Piscium), Lambda (λ Piscium), TX Piscium, Iota (ι Piscium), and Theta (θ Piscium).
Dubhe and Merak (Alpha and Beta Ursae Majoris), the two stars at the end of the bowl of the Big Dipper are often called the Pointers: a line from β to α and continued for about five times the distance between them arrives at the North Celestial Pole and the star Polaris (α UMi/Alpha Ursae Minoris), the North Star.
Rigil Kentaurus (α Centauri) and Hadar (β Centauri) are the Southern Pointers leading to the Southern Cross and thus helping to distinguish Crux from the False Cross.
Asterisms across multiple constellations
Other asterisms that are formed from stars in more than one constellation.
The Egyptian X is a large asterism which, like the Diamond of Virgo, is composed of a pair of equilateral triangles. Sirius (α CMa), Procyon (α CMi), and Betelgeuse (α Ori) form one to the North (Winter Triangle) while Sirius, Naos (ζ Pup), and Phakt (α Col) form another to the South. Unlike the Diamond, however, these triangles meet, not base-to-base, but vertex-to-vertex. The name derives from both the shape and, because the stars straddle the Celestial Equator, it is more easily seen from south of the Mediterranean than in Europe.
The Lozenge is a small diamond formed from three stars – Eltanin, Grumium, and Rastaban (Gamma, Xi, and Beta Draconis) – in the head of Draco and one – Iota Herculis – in the foot of Hercules.
The diamond-shaped False Cross is composed of the four stars Alsephina (δ Velorum), Markeb (κ Velorum), Avior (ε Carinae), and Aspidiske (ι Carinae). Although its component stars are not quite as bright as those of the Southern Cross, it is somewhat larger and better shaped than the Southern Cross, for which it is sometimes mistaken, causing errors in astronavigation. Like the Southern Cross, three of its main four stars are whitish and one orange.
The Northern Y is formed by four prominent stars, Arcturus (α Boötis), Seginus (γ Boötis), Alphecca (α Coronae Borealis), and centered on Izar (ε Boötis). From the United Kingdom in particular, where there is serious light pollution in many areas and also twilight much of the night when these constellations appear, this "Y" is often visible while other stars of Boötes and Corona Borealis are not.
The Lightning Bolt, aligned north to south, consists of the stars Epsilon Pegasi, Alpha Aquarii, Beta Aquarii and Delta Capricorni. Easily visible to naked eyes even in light polluted skies, the asterism is useful for orienting among three constellations.
The Serpent Bowl is a large curved asterism spanning 3.5 hours of right ascension, from mid-northern latitudes best seen in July and August evenings. From west to east, it includes the stars Delta, Alpha and Epsilon Serpentis, Delta, Epsilon, Upsilon, Zeta and Eta Ophiuchi, Xi Serpentis, Nu and Tau Ophiuchi, Eta and Theta Serpentis.
The Eagle Tail Corona is a flattened curved figure in the tail of Aquila and extending into Scutum. It consists of the stars 14, 15, Lambda and 12 Aquilae, Eta Scuti, HD 174208, R and Beta Scuti. The compact open cluster Messier 11 is also aligned with the curve.
Telescopic asterisms
Asterisms range from the large and obvious to the small, and even telescopic.
The 37 or LE of NGC 2169, in Orion.
The Engagement Ring in Ursa Minor has the north star Polaris as the diamond, at one end of a ring of much fainter stars about one degree across.
The Broken Engagement Ring in Ursa Major at 10:51 / +56°10' (preceding β Ursae Majoris, Merak).
The Christmas Tree shape of the Christmas Tree Cluster, in Monoceros. It is made up of about approximately 40 stars.
The Coathanger, in Vulpecula, also known as Brocchi's Cluster (see image at top).
Kemble's Cascade, a chain of stars that ends in open cluster NGC 1502, in Camelopardalis.
Napoleon's Hat (Picot 1), in Bootes (south of α Bootis, Arcturus).
The Ring of the Nibelungen (Ferrero 27) in Draco, named after the 1857 German epic drama, at 15:57 / +62°32' (near galaxy NGC 6015).
The V-shaped Messier 73 in Aquarius, determined to be an asterism in 2002.
| Physical sciences | Celestial sphere: General | Astronomy |
1028552 | https://en.wikipedia.org/wiki/Nelumbo | Nelumbo | Nelumbo is a genus of aquatic plants with large, showy flowers. Members are commonly called lotus, though the name is also applied to various other plants and plant groups, including the unrelated genus Lotus. Members outwardly resemble those in the family Nymphaeaceae ("water lilies"), but Nelumbo is actually very distant from that family.
Nelumbo is an ancient genus, with dozens of species known from fossil remains since the Early Cretaceous. However, there are only two known living species of lotus. One is the better-known Nelumbo nucifera, which is native to East Asia, South Asia, Southeast Asia, and probably Australia and is commonly cultivated for consumption and use in traditional Chinese medicine. The other lotus is Nelumbo lutea, which is native to North America and the Caribbean. Horticultural hybrids have been produced between these two allopatric species.
Description
Ultrahydrophobicity
The leaves of Nelumbo are highly water-repellent (i.e. they exhibit ultrahydrophobicity) and have given the name to what is called the lotus effect. Ultrahydrophobicity involves two criteria: a very high water contact angle between the droplet of water and the leaf surface, and a very low roll-off angle. This means that the water must contact the leaf surface at exactly one, minuscule point, and any manipulation of the leaf by changing its angle will result in the water droplet rolling off of the leaf. Ultrahydrophobicity is conferred by the usually dense layer of papillae on the surface of the Nelumbo leaves, and the small, robust, waxy tubules that protrude off each papilla. This helps reduce the area of contact between the water droplet and the leaf.
Ultrahydrophobicity is said to confer a very important evolutionary advantage. As an aquatic plant with leaves that rest on the water's surface, the genus Nelumbo is characterized by its concentration of stomata on the upper epidermis of its leaves, unlike most other plants which concentrate their stomata on the lower epidermis, underneath the leaf. The collection of water on the upper epidermis, whether that be by rain, mist, or the nearby disturbance of water, is very detrimental to the leaf's ability to perform gas exchange through its stomata. Thus, Nelumbo's ultrahydrophobicity allows the water droplets to accumulate together very quickly, and then roll off of the leaf very easily at the slightest disturbance of the leaf, a process which allows its stomata to function normally without restriction due to blockage by water droplets.
Thermoregulation
An uncommon property of the genus Nelumbo is that it can generate heat, which it does by using the alternative oxidase pathway (AOX). This pathway involves a different, alternative exchange of electrons from the usual pathway that electrons follow when generating energy in mitochondria, known as the AOX, or alternative oxidase pathway.
The typical pathway in plant mitochondria involves cytochrome complexes. The pathway used to generate heat in Nelumbo involves cyanide-resistant alternative oxidase, which is a different electron acceptor than the usual cytochrome complexes. The plant also reduces ubiquitin concentrations while in thermogenesis, which allows the AOX in the plant to function without degradation. Thermogenesis is restricted to the receptacle, stamen, and petals of the flower, but each of these parts produce heat independently without relying on the heat production in other parts of the flower.
There are several theories about the function of thermogenesis, especially in an aquatic genus such as Nelumbo. The most common theory posits that thermogenesis in flowers attracts pollinators, for a variety of reasons. Heated flowers may attract insect pollinators. As the pollinators warm themselves while resting inside the flower, they deposit and pick up pollen onto and from the flower. The thermogenic environment might also be conducive to pollinator mating; pollinators may require a certain temperature for reproduction. By providing an ideal thermogenic environment, the flower is pollinated by mating pollinators. Others theorize that heat production facilitates the release of volatile compounds into the air to attract pollinators flying over water, or that the heat is recognizable in the dark by thermo-sensitive pollinators. None have been conclusively proven to be more plausible than the others.
After anthesis, the receptacle of the lotus transitions from a primarily thermogenic to a photosynthetic structure, as seen in the rapid and dramatic increase in photosystems, photosynthetically involved pigments, electron transport rates, and the presence of 13C in the receptacle and petals, all of which assist in increasing photosynthesis rates. After this transition, all thermogenesis in the flower is lost. Pollinators do not need to be attracted once the ovary is fertilized, and thus the receptacle's resources are better used when it is photosynthesizing to produce carbohydrates that can increase plant biomass or fruit mass.
Other plants utilize thermoregulation in their life cycles. Among these is the eastern skunk cabbage, which heats itself to melt any ice above it, and push through the ground in early spring. Also, the elephant yam, which heats its flowers to attract pollinators. In addition, the carrion flower, which heats itself to disperse water vapor through the air, carrying its scent further, thus attracting more pollinators.
Similar species
The leaves of Nelumbo can be distinguished from those of genera in the family Nymphaeaceae as they are peltate, that is they have fully circular leaves. Nymphaea, on the other hand, has a single characteristic notch from the edge in to the center of the lily pad. The seedpod of Nelumbo is very distinctive.
Taxonomy
Taxonomic history
The Cronquist system of 1981 recognizes the family Nelumbonaceae but places it in the water lily order Nymphaeales. The Dahlgren system of 1985 and Thorne system of 1992 both recognize the family and place it in its own order, Nelumbonales. The United States Department of Agriculture still classifies the lotus family within the water lily order.
There is residual disagreement over which family the genus should be placed in. Traditional classification systems recognized Nelumbo as part of the Nymphaeaceae, but traditional taxonomists were likely misled by convergent evolution associated with an evolutionary shift from a terrestrial to an aquatic lifestyle. In the older classification systems it was recognized under the order Nymphaeales or Nelumbonales.
Modern classification
Nelumbo is currently recognized as the only living genus in Nelumbonaceae, one of several distinctive families in the eudicot order of the Proteales. Its closest living relatives, the (Proteaceae and Platanaceae), are shrubs or trees.
The APG IV system of 2016 recognizes Nelumbonaceae as a distinct family and places it in the order Proteales in the eudicot clade, as do the earlier APG III and APG II systems.
Phylogeny
There are several fossil species known from Cretaceous, Paleogene and Neogene aged strata throughout Eurasia and North America. Despite the ancient origins of this genus and the wide geographic separation of the two extant species (N. nucifera and N. lutea), phylogenetic evidence indicates that they diverged rather recently, during the early Pleistocene (about 2 million years ago).
Species
Extant species
Nelumbo lutea Willd. – American lotus (Eastern United States, Mexico, Greater Antilles, Honduras)
Nelumbo nucifera Gaertn. – sacred or Indian lotus, also known as the Rose of India and the sacred water lily of Hinduism and Buddhism. It is the national flower of India and Vietnam. Its roots and seeds are also used widely in cooking in East Asia, South Asia and Southeast Asia.
Fossil species
Nearly 30 fossil species are known from the mid-Cretaceous to the present.
†Nelumbo aureavallis Hickey – Eocene (North Dakota), described from leaves found in the Golden Valley Formation in North Dakota, USA.
†Nelumbo changchangensis Eocene, (Hainan Island, China), described from several fossils of leaves, seedpods, and rhizomes from the Eocene-aged strata in the Changchang Basin, of Hainan Island.
†Nelumbo choffati Early Cretaceous (Portugal), leaves known from the Albian. One of the earliest known species.
†Nelumbo jiayinensis Late Cretaceous, (Heilongjiang, China), leaves described from the Santonian-aged Yong'ancun Formation
†Nelumbo lusitanica Early Cretaceous (Portugal), leaves known from the Albian. One of the earliest known species.
†Nelumbo minima Pliocene (Netherlands), described from leaves and seedpods that suggest a very small plant. Originally described as a member of the genus Nelumbites, as "Nelumbites minimus."
†Nelumbo nipponica Eocene-Miocene, fossil leaves are known from Eocene-aged strata in Japan, and Miocene-aged strata in Russia.
†Nelumbo orientalis Cretaceous (Japan), fossils found in Cretaceous-aged strata of Japan. The Sarao Formation, which they are known from, was formerly considered of Early Cretaceous age, but more recent studies support a Maastrichtian age for it.
†Nelumbo protolutea Eocene (Mississippi), fossils of leaves strongly suggest a plant similar in form to the American lotus.
†Nelumbo weymouthi Early Cretaceous (Wyoming, US), leaves known from the Albian. One of the earliest known species.
Etymology
The genus name is derived from neḷum, the name for Nelumbo nucifera.
Uses
The entire plant can be eaten either raw or cooked. The underwater portion is high in starch. The fleshy parts can be dug from the mud and baked or boiled. The young leaves can be boiled. The seeds are palatable and can be eaten raw or dried and ground into flour. The stem fibers are also used to make lotus silk.
Culture
The sacred lotus, N. nucifera, is sacred in both Hinduism and Buddhism. It is the floral emblem of both India and Vietnam.
| Biology and health sciences | Others | null |
1028614 | https://en.wikipedia.org/wiki/Nelumbo%20nucifera | Nelumbo nucifera | Nelumbo nucifera, also known as sacred lotus, Indian lotus, or simply lotus, is one of two extant species of aquatic plant in the family Nelumbonaceae. It is sometimes colloquially called a water lily, though this more often refers to members of the family Nymphaeaceae.
Lotus plants are adapted to grow in the flood plains of slow-moving rivers and delta areas. Stands of lotus drop hundreds of thousands of seeds every year to the bottom of the pond. While some sprout immediately and most are eaten by wildlife, the remaining seeds can remain dormant for an extensive period of time as the pond silts in and dries out. During flood conditions, sediments containing these seeds are broken open, and the dormant seeds rehydrate and begin a new lotus colony.
Under favorable circumstances, the seeds of this aquatic perennial may remain viable for many years, with the oldest recorded lotus germination being from seeds 1,300 years old recovered from a dry lakebed in northeastern China. Therefore, the Chinese regard the plant as a symbol of longevity.
It has a very wide native distribution, ranging from central and northern India (at altitudes up to in the southern Himalayas), through northern Indochina and East Asia (north to the Amur region; the Russian populations have sometimes been referred to as Nelumbo komarovii, with isolated locations at the Caspian Sea. Today, the species also occurs in southern India, Sri Lanka, virtually all of Southeast Asia, New Guinea, and northern and eastern Australia, but this is probably the result of human translocations. It has a very long history ( 3,000 years) of being cultivated for its edible seeds and is commonly cultivated in water gardens. It is the national flower of India and Vietnam.
Classification
The lotus is often confused with the true water lilies of the genus Nymphaea, in particular N. caerulea, the "blue lotus." In fact, several older systems, such as the Bentham & Hooker system (which is widely used in the Indian subcontinent), refer to the lotus by its old synonym, Nymphaea nelumbo.
While all modern plant taxonomy systems agree that this species belongs in the genus Nelumbo, the systems disagree as to which family Nelumbo should be placed in or whether the genus should belong in its own unique family and order. According to the APG IV system, N. nucifera, N. lutea, and their extinct relatives belong in Proteales with the protea flowers due to genetic comparisons. Older systems, such as the Cronquist system, place N. nucifera and its relatives in the order Nymphaeles based on anatomical similarities. According to the APG IV classification, the closest relatives of Nelumbo include the sycamores (Platanaceae).
Botany
The lotus roots are planted in pond or river bottom soil, while the leaves float on the water's surface or are held well above it. The leaf stalks (petioles) can be up to long, allowing the plant to grow in water to that depth. The peltate leaf blade or lamina can have a horizontal spread of . The leaves may be as large as in diameter.
Flower
The flowers are usually found on thick stems rising several centimeters above the leaves. They are showy and grow up to in diameter.
Some cultivated varieties have extraordinary numbers of petals. For example, the Chinese variety qian ban lian ("thousand petals lotus") can have between 3000 and 4000 petals in a single blossom and the Japanese variety ohmi myoren ("strange lotus") can have between 2000 and 5000 petals, the greatest number recorded for any species of plant.
Researchers report that the lotus has the remarkable ability to regulate the temperature of its flowers within a narrow range, just as humans and other warm-blooded animals do. Roger S. Seymour and Paul Schultze-Motel, physiologists at the University of Adelaide in Australia, found that lotus flowers blooming in the Adelaide Botanic Gardens maintained a temperature of , even when the air temperature dropped to . They suspect the flowers may be doing this to attract cold-blooded insect pollinators. Studies published in the journals Nature and Philosophical Transactions: Biological Sciences in 1996 and 1998 were important contributions in the field of thermoregulation in plants. Two other species known to be able to regulate their temperature include Symplocarpus foetidus and Thaumatophyllum bipinnatifidum. The red tiger lotus is native to West Africa, including Nigeria and Cameroon, and thrives in slow-moving water.
Seed
A fertilized lotus flower bears fruit that contains a cluster of 10 to 30 seeds. Each seed is ovoid 1–2.5 cm wide by 1–1.5 cm long with a brownish coat. Lotus seeds can remain viable after long periods of dormancy. In 1994, a seed from a sacred lotus, dated at roughly 1,300 years old ± 270 years, was successfully germinated.
The traditional sacred lotus is only distantly related to Nymphaea caerulea, but possesses similar chemistry. Both Nymphaea caerulea and Nelumbo nucifera contain the alkaloids nuciferine and aporphine.
The genome of the sacred lotus was sequenced in May 2013. A dedicated genome database lists additional genome assemblies sequenced since then.
Cultivation
The sacred lotus grows in water up to deep. The minimum water depth is about . In colder climates, having a deeper water level protects the tubers more effectively, and overall is helpful for better growth and flowering. The sacred lotus germinates at temperatures above . Most varieties are not naturally cold-hardy, but may readily adapt to living outdoors year-round in USDA hardiness zones 6 through 11 (with some growers having success in zones as low as 4 or 5); the higher the zone's number, the greater the adaptability of the plants. In the growing season, from April to September (in the northern hemisphere), the average daytime temperature needed is . In regions with low light levels in winter, the sacred lotus has a period of dormancy. The tubers are not cold-resistant, if removed from water, and exposed to the air; when kept underwater in soil, the energy-rich tubers can overwinter temperatures below . If the plants are taken out of the water for wintertime storage (mostly in exceptionally cold climates), the tubers and roots must be stored in a stable, frost-free location, such as a garage, preferably in a cardboard box or container filled completely with vermiculite or perlite. Care must be taken to fully insulate the tubers.
Planting
The sacred lotus requires a nutrient-rich and loamy soil. In the beginning of the summer period (from March until May in the northern hemisphere), a small part of rhizome with at least one eye is either planted in ponds or directly into a flooded field. There are several other propagation ways via seeds or buds. Furthermore, tissue culture is a promising propagation method for the future to produce high volumes of uniform, true-to-type, disease-free materials.
The first step of the cultivation is to plough the dry field. One round of manure is applied after ten days, before flooding the field. To support a quick initial growth, the water level is relatively low and increases when plants grow. Then a maximum of approximately with grid spacing of are used to plant directly into the mud below the soil surface.
Harvest
The stolon is ready to harvest two to three months after planting. It must be harvested before flowering. Harvesting the stolon is done by manual labor. For this step, the field is not drained. The stolon is pulled out of the water by pulling and shaking the young leaves in the shallow water.
The first leaves and flowers can be harvested three months after planting. Flowers can be picked every two days during summer and every three days during the colder season. Four months after planting, the production of flowers has its climax. The harvest of flowers is usually done by hand for three to four months.
Seeds and seed pods can be harvested when they turn black four to eight months after planting. After sun drying for two to three days, they are processed by mechanical tools to separate seed coats and embryos.
The rhizomes mature to a suitable stage for eating in approximately six to nine months. Early varieties are harvested in July until September and late varieties from October until March, after the ponds or fields are drained. The large, starch-rich rhizomes are easy to dig out of the drained soil. In small-scale production, they are harvested by hand using fork-like tools. In Japan and on bigger farms, manual labour harvesting is fully replaced by machines.
Varieties and cultivars
Lotus varieties have been classified according to their use into three types: rhizome lotus, seed lotus, and flower lotus. Varieties that show more than one of these characteristics are classified by the strongest feature. Regarding production area in China, rhizome lotus has the largest area with , followed by seed lotus with .
Rhizome lotus
Rhizome lotus cultivars produce a higher yield and higher quality rhizomes than seed or flower lotus cultivars. Furthermore, this group grows tall and produces few to no flowers.
Cultivars can be classified by harvest time or by the depth of rhizomes into these types:
Pre-mature (early) cultivars are harvested before the end of July, serotinous (late) cultivars from September on, and mid-serotinous or mid-matutinal cultivars are in between these harvest times. Using pre-mature cultivars, rhizomes can be harvested earlier and sold for a higher price.
Adlittoral, deep, and intermediate cultivars are distinguished according to the depth in which the rhizomes grow underground. Adlittoral cultivars range from depth and are often premature. They develop faster due to higher temperatures in surface soil layers. When harvested in July, adlittorals have higher yields than deeper-growing cultivars, but not necessarily when harvested in September. Rhizomes of adlittoral cultivars are crisp and good for frying purposes. Deep cultivars grow more than deep. They are often serotinous and can harvest high yields. Their rhizomes are starch-rich.
The main popular Nelumbo nucifera cultivars in China are Elian 1, Elian 4, Elian 5, 9217, Xin 1, and 00–01. The average yield of these cultivars is 7.5–15 t/ha (3.3–6.7 tons/acre) of harvest in July and 30–45 t/ha (13–20 tons/acre) of harvest in September. In Australia, the cultivar grown for the fresh rhizome market in Guangdong and Japan, the common rhizome cultivars are Tenno and Bitchu.
Seed lotus
The characteristics of seed lotus cultivars are a large number of carpels and seed sets as well as large seeds with better nutritional properties. Roots of these varieties are thin, fibrous, and do not form good rhizomes. The main popular cultivars for seed production in China are Cunsanlian, Xianglian 1, Zilian 2, Jianlian, Ganlian 62, and Taikong 36. The average yield of these cultivars in China is 1.05–1.9 t/ha (0.5–0.8 tons/acre) of dry seeds and weight of thousand seeds between . Green Jade and Vietnam-Red are recommended cultivars for seed production in Australia.
Flower lotus
Flower lotus cultivars are used exclusively for ornamental purpose, producing many flowers and the lowest plant height.
The seed production of flower lotus is typically poor regarding yield and quality. Flower types differ in the number of petals (single petals, double petals, or multi-petals) and their colours range from single colour in white, yellow, pink, and red to bi-colour, most often of white petals with pink tips or highlights.
The flowers are capable of producing ink used by artists such as Morrison Polkinghorne to produce abstract images of the landscapes of southern Asia.
One example of a flower lotus is Wanlian. Also known as bowl lotus, wanlians are any miniature cultivars of N. nucifera sized between . Bowl lotuses come in various colours and numbers of petals, and they bloom longer than other species of lotus. But together with the rhizome, their seeds are often too small or too hard to be edible.
The sacred lotus may be crossed with the yellow lotus to produce interspecific hybrids. A few varieties have been produced with differing appearances.
Farming
About 70% of lotus for human consumption is produced in China. In 2005, the cultivation area in China was estimated at . A majority of lotus production takes place in managed farming systems in ponds or flooded fields like rice.
The most widely used system is crop rotation with rice and vegetables. This system is applicable if the propagule (small piece of rhizome) can be planted early in the year. The rhizomes are harvested in July, after which rice can be planted in the same field. Rice is then harvested in October. From November until March, the field stays either free or terricolous vegetables, such as cabbage or spinach, are planted. Alternatively, the vegetable can also be planted after the harvest of the lotus.
Another alternative way is not to harvest the lotus rhizome, even though it is already ripe. A terricolous vegetable is planted between the rhizomes into the drained field. The rhizomes are then harvested next March.
A third way is to plant lotus in ponds or fields and raise aquatic animals such as fish, shrimp, or crab in the same field. A more efficient use of the water for both, the aquatic animals and lotus production has been identified with this planting pattern.
Use
Religious
Lotus flowers are widely used as offerings to most female deities, especially Lakshmi, in Hindu temples. Among male deities, lotuses are offered to Vishnu for prosperity and to Shiva for salvation. Garlands made of lotuses are used for adorning deities and lotus petals are used in puja. Lotus seeds are also used in prayer beads. Lotuses are also offered to the Buddha in most Buddhist temples. Lotus is also widely used in Varamala (hindu wedding garland).
Political
Lotus is the national flower of the Republic of India. Bharatiya Janata Party (BJP), the largest political party in the world with around 200 million active members, uses lotus as its party symbol. BJP, in the past, have used lotuses in multiple colors - pink,white,blue,red and saffron, in their party flag. In recent times, as a part of branding strategy and to make it easier for voters, BJP started using lotus logo in black-and-white in most settings as Electronic Voting Machines (EVM) in India only allow black-and-white logos.
Culinary
Rhizomes
The rhizomes of lotus (, , , , Sindhi Beeh, ) are consumed as a vegetable in Asian countries, extensively in China, Japan, India, Pakistan (Sindh), sold whole or in cut pieces, fresh, frozen, or canned. They are fried or cooked mostly in soups, soaked in syrup or pickled in vinegar (with sugar, chili and garlic). Lotus rhizomes have a crunchy texture and are a classic dish at many banquets, where they are deep-fried, stir-fried, or stuffed with meats or preserved fruits. Salads with prawns, sesame oil or coriander leaves are also popular. Fresh lotus root slices are limited by a fast browning rate. Lotus root tea is consumed in Korea.
Lotus root is a popular vegetable in Sri Lanka, where it is often cooked in coconut milk gravy. In India, lotus root (also known as kamala kakaṛī in Hindi) is cooked as a dry curry or Sabzî.
Japan is one of the primary users of the rhizomes, representing about 1% of all vegetables consumed. Japan grows its own lotus but still must import 18,000 tons of lotus rhizome each year, of which China provides 15,000 tons yearly.
Rhizomes contain high amounts of starch (31.2%) without characteristic taste or odor. The texture is comparable to a raw potato. The binding and disintegration properties of isolated Nelumbo starch have been compared with maize and potato starch; Nelumbo starch is shown to be superior as an adjuvant in the preparation of tablets. When dried, N. nucifera is also made into flour, another popular use of this vegetable.
Pips
Lotus pip tea is consumed in Korea.
Seeds
Fresh lotus seeds () are nutritious but also vulnerable to microbial contamination, especially fungal infections. Therefore, mostly dry lotus seed-based products are found on the market. Traditional sun baking combined with charcoal processing dries the seeds but results in a loss of nutrients. Freeze-dried lotus seeds have a longer shelf life and maintain original nutrients, while no differences in flavour are found after rehydration compared to fresh lotus seeds.
Dry stored lotus seeds are sensitive to moisture and mold infestation; researchers continue to explore new ways to preserve fresh lotus seeds, such as radiation processing.
Lotus seeds can be processed into fillings for moon cake, lotus seed noodles and food in the forms of paste, fermented milk, rice wine, ice cream, popcorn (phool makhana), and others, with lotus seeds as the main raw material. Traditional Chinese medicine claims that fresh lotus seed wine has thirst-quenching, spleen-healing, and anti-diarrheal advantages after drinking, attributed to unspecified bioactive compounds. Lotus seed tea is consumed in Korea, and lotus embryo tea is consumed in China and Vietnam.
Stems
Young lotus stems are used as a salad ingredient in Vietnamese cuisine and as a vegetable ingredient for some soup and curry in Thailand, such as keang som sai bua (, lotus stem sour soup) and keang kati sai bua (, lotus stem in coconut milk curry).
In northern and eastern regions of India, the stalk of the flower is used to prepare a soup, kamala gaṭṭē kī sabzī () and an appetizer, kamala kakaṛī pakauṛē (). In South Indian states, the lotus stem is sliced, marinated with salt to dry, and the dried slices are fried and used as a side dish. In Kerala () and Tamil Nadu, this end product is called thamara vathal.
In the Philippines, an indigenous variety called tukal is used as the main ingredient in dishes with coconut milk. The stems and petals can be bought in markets when in season.
Leaves
In China and Korea, lotus leaf tea () is made from the leaves of the lotus. It is also used as a wrap for steaming rice and sticky rice and other steamed dishes in Southeast Asian cuisine, such as lo mai gai in Chinese cuisine or kao hor bai bua (), fried rice wrapped in lotus leaf in Thai cuisine.
Vietnamese also use lotus leaves to wrap green young rice, cốm, which is eaten in autumn. The leaves impart a unique scent to the soft, moist rice.
Flowers
In Korea, lotus flower tea () is made from the dried petals of the white lotus.
The stamens can be dried and made into a fragrant herbal tea (), or used to impart a scent to tea leaves (particularly in Vietnam). This Vietnamese lotus tea is called trà sen, chè sen, or chè ướp sen.
Risks
The petals, leaves, and rhizome can also all be eaten raw, but there is a risk of parasite transmission (e.g., Fasciolopsis buski): it is therefore recommended that they be cooked before eating.
Use in water treatment
Nelumbo nucifera shows high potential for usage in wastewater treatment removing polluting compounds and heavy metals. It is able to grow in variable water conditions and in low light intensity. Various studies show the successful use of N. nucifera to counteract water eutrophication. The leaves of the floating lotus reduce sunlight reaching the lower part of the water. This suppresses algae growth in N. nucifera aquatic systems and thus, the oxygen content is up to 20% higher than in other aquatic plant systems. Due to intense agricultural practices, nitrogen and phosphorus pollution are major problems in aquatic systems. N. nucifera is able to assimilate a higher content of phosphorus than aquatic plants currently used for water remediation (such as water hyacinth). It also assimilates nitrogen ("denitrification") and creates a habitat for bacterial growth in the water body. Through rhizofiltration, heavy metals – including arsenic, copper, and cadmium – can be removed efficiently from the water. The results observed are impressive showing 96% of copper and 85% cadmium metals removed after a seven-day incubation period. The accumulation of heavy metals doesn't show morphological symptoms of metal toxicity; however, the rhizome quality for human consumption needs further study.
Storage and commercialization
Currently, most rhizomes are consumed fresh, and it is not common to store them due to their poor shelf life performance.{ This limits export possibilities for low-income production countries in Asia. Rhizomes quickly lose water, oxidation occurs and nutrient composition changes within a short time after harvest. Optimal storage temperatures range between . There are three different approaches to storing rhizomes. By stacking the rhizomes, they are storable and remain fresh for about three weeks. Special stacking with silver sand and soil results in five to six layers that prevent water loss, thus, the rhizome stays fresh for up to two months. However, the method is not suitable for commercialization but rather for home use. Hydrogen sulfide fumigation reduces enzymatic browning and therefore ensures rhizome quality. Dipping the rhizomes in a salt solution prevents oxidation and bacterial reproduction, which allows storage for up to five months and greater export ability. This treatment is related to the high cost and inefficient cleaning process before eating the rhizomes.
Use in bioengineering
Nelumbo nucifera contains some thermal-stable proteins that might be useful in protein bioengineering processes. The proteins are characterized by seed longevity used for cell protection and repair under stress. There are also several indications that compounds of N. nucifera are used in drug fabrication in human health research for multiple purposes. Lotus leaves possess hydrophobic characteristics, attributed to a waxy coat that prevents water from adhering to the surface. This attribute has influenced the conception of the "lotus effect" in biomimicry and engineering, guiding the design of materials that resist water and remain self-cleaning. Researchers at the National University of Singapore have utilized the water-repelling structure as inspiration for developing eAir, an aero-elastic sensor capable of detecting subtle pressure changes or other environmental stimuli.
Other uses
The distinctive dried seed heads, which resemble the spouts of watering cans, are widely sold throughout the world for decorative purposes and for dried flower arranging.
In Asia, the petals are sometimes used for garnish, while the large leaves are used as a wrap for food, not frequently eaten (for example, as a wrapper for zongzi). Lotus leaves are also used to serve food in various cultures.
A unique fabric called lotus silk, from the lotus plant fibers, is produced only at Inle Lake, Myanmar, and in Siem Reap, Cambodia. This thread is used for weaving special robes for Buddha images called kya thingan (lotus robe).
Chemical composition
The flavonol miquelianin, as well as the alkaloids (+)-(1R)-coclaurine and (−)-(1S)-norcoclaurine, can be found in the leaves of N. nucifera. The plant also contains nuciferine, neferine, and many other benzylisoquinoline alkaloids with medicinal properties.
Health properties and nutrients
Traditional medicine
All parts of Nelumbo nucifera are edible, with the rhizome and seeds being the main consumption parts. Traditionally rhizomes, leaves, and seeds have been used as folk medicines, Ayurveda, Chinese traditional medicine, and oriental medicine. In Chinese medicine, seeds are still used as ().
Lotus rhizomes and seeds and their processing by-products are widely consumed in Asia, the Americas, and Oceania for high content of physiologically active substances. Especially in China, lotus seeds are popular with a cultural history going back about 3000 years. As early as the Han Dynasty, lotus seeds were already recorded as sweet, astringent, nourishing the heart and kidney in "Shen Nong's Herbal Classic". Nowadays there are 22 varieties for the four known Chinese lines, which are found particularly in Jianning (still called "the town of Jianlian lotus seeds in China") and Guangchang ("the town of white lotus seeds in China").
These days the perennial aquatic herb is gaining popularity because of its nutraceutical and historical importance It will be of economic value if the different parts of lotus can be developed as functional food. Because of its special role in human health and richness in nutrients and bioactive substances, the Chinese Ministry of Health approved the use of N. nucifera as both "food and medicine."
Rhizomes
The rhizomes are long, in diameter, yellowish white to yellowish brown, smooth, and with nodes and internodes.
Lotus root is a moderate-calorie root vegetable (100 g of root stem provides about 74 calories) and is composed of several vitamins, minerals, and nutrients: 83.80% water, 0.11% fat, 1.56% reducing sugar, 0.41% sucrose, 2.70% crude protein, 9.25% starch, 0.80% fiber, 0.10% ash and 0.06% calcium. 100 g of root provides 44 mg of vitamin C or 73% of daily recommended values (RDA).
Lotus rhizome and its extracts have shown diuretic, psychopharmacological, anti-diabetic, anti-obesity, hypoglycemic, antipyretic and antioxidant activities.
Seeds
Lotus seeds are mostly oval or spherical, with sizes varying according to varieties. They are generally 1.2–1.8 cm long, with diameters ranging from 0.8 to 1.4 cm and a weight of 1.1–1.4 g. After lotus seeds have been decorticated and peeled, they are edible and rich in nutrients and can be dried for storage. Their nutritional values can differ due to culture environments and varieties.
Not only do these seeds contain proteins of high quality and are rich in a variety of essential amino acids including high contents of albumin (42%) and globulin (27%), they also contain unsaturated fatty acids, carbohydrates, vitamins, calcium, iron, zinc, phosphorus and other trace elements. They also provide water-soluble polysaccharides, alkaloids, flavonoids, superoxide dismutase, and other bioactive components.
Lotus seeds also contain particularly large amounts of vitamins, including VB1, VB2, VB6 and Vitamin E.
The functional components (polyphenols, protein, polysaccharides) in N. nucifera seeds can help combat high blood pressure, diabetes, and gallstones.
After lotus seed germination, crude protein and fat levels in the endosperm significantly increase. It is therefore an important method to enhance its nutritional quality.
Cultural and religious significance
Nelumbo nucifera is a lotus species with historical, cultural and spiritual significance. It is a sacred flower in both Hinduism and Buddhism, representing the path to spiritual awakening and enlightenment. Lotus is closely associated with Goddess Lakshmi and her consort Vishnu. In Hindu iconography, Lakshmi is shown either in seated posture or as standing on top of a lotus flower and holding lotuses in two of her four hands. Vishnu is always portrayed with a lotus in one of his hands. Lotus or Padma exemplifies the path of one who leads a dharmic life. One can stay in the darkness, trapped in the comfort of the decaying leaves and murky water or struggle to lead a dharmic life upholding righteousness, following the light and eventually liberate oneself of his/her past thereby leading a prosperous peaceful happy life. Lotus also signifies moksha. The roots of lotus plant intertwined in mud representing the cycle of life and death and the lotus flower symbolizing the serene blissful state of moksha. Lotus is also used to depict the awakening of Kundalini energy in the tantric system of Hinduism.
In Hindu literature, lotus petals are used as a simile to adore beautiful eyes of gods and goddesses. For instance, the word Kamalanayana (lotus-eyed) refers to Vishnu and it extols his beautiful eyes and praise him for being the one who holds the gaze/adoration of Kamala (Lakshmi).
Epic Mahabharata details the use of multi-tiered military formation that resembles a blooming lotus called Padmavyuha and its disc-shaped variant Chakravyuha in the 18-day long Kurukshethra war.
Padmasana (Lotus Pose) in Yoga and Sahasrara - thousand petalled lotus chakra in tantra are some examples of the use of lotus imagery.
Lotus has also inspired Indian/Bharatiya architecture and can be seen in domes resembling lotus buds, pillars resembling lotus stalk, arches resembling inverted lotus and other lotus engravings in Hindu/Buddhist temples.In Asian art, a lotus throne is a stylized lotus flower used as the seat or base for a figure. It is the normal pedestal for divine figures in Buddhist art and Hindu art and is often seen in Jain art. Originating in Indian art, it followed Indian religions to East Asia in particular. Lotus flowers are also often held by figures.
The Nelumbo nucifera, which is also called (Nilufar Abi in Persian), can be seen in many reliefs of the Achaemenid period (552 BC) such as the statue of Anahita in the Persepolis. The lotus flower was included in Kaveh the blacksmith's Derafsh and later as the flag of the Sasanian Empire Derafsh Kaviani. Today, it is the symbol of Iranians Solar Hijri calendar.
The lotus flower is also mentioned in the Bible. The lotus flower also holds cultural and religious significance in Ismaili Muslim and related South Asian traditions. For example, in South Asian Ismaili literature, the lotus is compared to the pure soul. A poem describes the lotus' beauty, describing how its delicate white petals remain pure and beautiful, despite its murky environment. Similarly, a pure soul is part of this world, yet is not of this world, much like the circumstances of the lotus. The poem further emphasizes the importance of true knowledge or gnosis, which is likened to the pure rainwater that allows the lotus to flourish. Ismaili belief holds that the true guide provides this true knowledge, without which the pure soul cannot survive. Just as the lotus flower would rather die than drink from a reeking swamp, the pure soul also seeks nourishment solely through true knowledge.
In Chinese culture, the lotus is known as “Lianhua” (蓮花). The Chinese characters “蓮” (蓮) and “花” (花) represent the leaves and flowers of the plant respectively. The lotus holds important cultural significance in Chinese Buddhism, symbolizing purity, enlightenment, and the unfolding of the spiritual self.
| Biology and health sciences | Proteales | Plants |
1029022 | https://en.wikipedia.org/wiki/Embryonic%20stem%20cell | Embryonic stem cell | Embryonic stem cells (ESCs) are pluripotent stem cells derived from the inner cell mass of a blastocyst, an early-stage pre-implantation embryo. Human embryos reach the blastocyst stage 4–5 days post fertilization, at which time they consist of 50–150 cells. Isolating the inner cell mass (embryoblast) using immunosurgery results in destruction of the blastocyst, a process which raises ethical issues, including whether or not embryos at the pre-implantation stage have the same moral considerations as embryos in the post-implantation stage of development.
Researchers are currently focusing heavily on the therapeutic potential of embryonic stem cells, with clinical use being the goal for many laboratories. Potential uses include the treatment of diabetes and heart disease. The cells are being studied to be used as clinical therapies, models of genetic disorders, and cellular/DNA repair. However, adverse effects in the research and clinical processes such as tumors and unwanted immune responses have also been reported.
Properties
Embryonic stem cells (ESCs), derived from the blastocyst stage of early mammalian embryos, are distinguished by their ability to differentiate into any embryonic cell type and by their ability to self-renew. It is these traits that makes them valuable in the scientific and medical fields. ESCs have a normal karyotype, maintain high telomerase activity, and exhibit remarkable long-term proliferative potential.
Pluripotent
Embryonic stem cells of the inner cell mass are pluripotent, meaning they are able to differentiate to generate primitive ectoderm, which ultimately differentiates during gastrulation into all derivatives of the three primary germ layers: ectoderm, endoderm, and mesoderm. These germ layers generate each of the more than 220 cell types in the adult human body. When provided with the appropriate signals, ESCs initially form precursor cells that in subsequently differentiate into the desired cell types. Pluripotency distinguishes embryonic stem cells from adult stem cells, which are multipotent and can only produce a limited number of cell types.
Self renewal and repair of structure
Under defined conditions, embryonic stem cells are capable of self-renewing indefinitely in an undifferentiated state. Self-renewal conditions must prevent the cells from clumping and maintain an environment that supports an unspecialized state. Typically this is done in the lab with media containing serum and leukemia inhibitory factor or serum-free media supplements with two inhibitory drugs ("2i"), the MEK inhibitor PD03259010 and GSK-3 inhibitor CHIR99021.
Growth
ESCs divide very frequently due to a shortened G1 phase in their cell cycle. Rapid cell division allows the cells to quickly grow in number, but not size, which is important for early embryo development. In ESCs, cyclin A and cyclin E proteins involved in the G1/S transition are always expressed at high levels. Cyclin-dependent kinases such as CDK2 that promote cell cycle progression are overactive, in part due to downregulation of their inhibitors. Retinoblastoma proteins that inhibit the transcription factor E2F until the cell is ready to enter S phase are hyperphosphorylated and inactivated in ESCs, leading to continual expression of proliferation genes. These changes result in accelerated cycles of cell division. Although high expression levels of pro-proliferative proteins and a shortened G1 phase have been linked to maintenance of pluripotency, ESCs grown in serum-free 2i conditions do express hypo-phosphorylated active Retinoblastoma proteins and have an elongated G1 phase. Despite this difference in the cell cycle when compared to ESCs grown in media containing serum these cells have similar pluripotent characteristics. Pluripotency factors Oct4 and Nanog play a role in transcriptionally regulating the embryonic stem cell cycle.
Uses
Due to their plasticity and potentially unlimited capacity for self-renewal, embryonic stem cell therapies have been proposed for regenerative medicine and tissue replacement after injury or disease. Pluripotent stem cells have shown promise in treating a number of varying conditions, including but not limited to: spinal cord injuries, age related macular degeneration, diabetes, neurodegenerative disorders (such as Parkinson's disease), AIDS, etc. In addition to their potential in regenerative medicine, embryonic stem cells provide a possible alternative source of tissue/organs which serves as a possible solution to the donor shortage dilemma. There are some ethical controversies surrounding this though (see Ethical debate section below). Aside from these uses, ESCs can also be used for research on early human development, certain genetic disease, and in vitro toxicology testing.
Utilizations
According to a 2002 article in PNAS, "Human embryonic stem cells have the potential to differentiate into various cell types, and, thus, may be useful as a source of cells for transplantation or tissue engineering."
Tissue engineering
In tissue engineering, the use of stem cells are known to be of importance. In order to successfully engineer a tissue, the cells used must be able to perform specific biological functions such as secretion of cytokines, signaling molecules, interacting with neighboring cells, and producing an extracellular matrix in the correct organization. Stem cells demonstrates these specific biological functions along with being able to self-renew and differentiate into one or more types of specialized cells. Embryonic stem cells is one of the sources that are being considered for the use of tissue engineering. The use of human embryonic stem cells have opened many new possibilities for tissue engineering, however, there are many hurdles that must be made before human embryonic stem cell can even be utilized. It is theorized that if embryonic stem cells can be altered to not evoke the immune response when implanted into the patient then this would be a revolutionary step in tissue engineering. Embryonic stem cells are not limited to tissue engineering.
Cell replacement therapies
Research has focused on differentiating ESCs into a variety of cell types for eventual use as cell replacement therapies. Some of the cell types that have or are currently being developed include cardiomyocytes, neurons, hepatocytes, bone marrow cells, islet cells and endothelial cells. However, the derivation of such cell types from ESCs is not without obstacles, therefore research has focused on overcoming these barriers. For example, studies are underway to differentiate ESCs into tissue specific cardiomyocytes and to eradicate their immature properties that distinguish them from adult cardiomyocytes.
Clinical potential
Researchers have differentiated ESCs into dopamine-producing cells with the hope that these neurons could be used in the treatment of Parkinson's disease.
ESCs have been differentiated to natural killer cells and bone tissue.
Studies involving ESCs are underway to provide an alternative treatment for diabetes. For example ESCs have been differentiated into insulin-producing cells, and researchers at Harvard University were able to produce large quantities of pancreatic beta cells from ESCs.
An article published in the European Heart Journal describes a translational process of generating human embryonic stem cell-derived cardiac progenitor cells to be used in clinical trials of patients with severe heart failure.
Drug discovery
Besides becoming an important alternative to organ transplants, ESCs are also being used in the field of toxicology, and as cellular screens to uncover new chemical entities that can be developed as small-molecule drugs. Studies have shown that cardiomyocytes derived from ESCs are validated in vitro models to test drug responses and predict toxicity profiles. ESC derived cardiomyocytes have been shown to respond to pharmacological stimuli and hence can be used to assess cardiotoxicity such as torsades de pointes.
ESC-derived hepatocytes are also useful models that could be used in the preclinical stages of drug discovery. However, the development of hepatocytes from ESCs has proven to be challenging and this hinders the ability to test drug metabolism. Therefore, research has focused on establishing fully functional ESC-derived hepatocytes with stable phase I and II enzyme activity.
Models of genetic disorder
Several new studies have started to address the concept of modeling genetic disorders with embryonic stem cells. Either by genetically manipulating the cells, or more recently, by deriving diseased cell lines identified by prenatal genetic diagnosis (PGD), modeling genetic disorders is something that has been accomplished with stem cells. This approach may very well prove valuable at studying disorders such as Fragile-X syndrome, Cystic fibrosis, and other genetic maladies that have no reliable model system.
Yury Verlinsky, a Russian-American medical researcher who specialized in embryo and cellular genetics (genetic cytology), developed prenatal diagnosis testing methods to determine genetic and chromosomal disorders a month and a half earlier than standard amniocentesis. The techniques are now used by many pregnant women and prospective parents, especially couples who have a history of genetic abnormalities or where the woman is over the age of 35 (when the risk of genetically related disorders is higher). In addition, by allowing parents to select an embryo without genetic disorders, they have the potential of saving the lives of siblings that already had similar disorders and diseases using cells from the disease free offspring.
Repair of DNA damage
Differentiated somatic cells and ES cells use different strategies for dealing with DNA damage. For instance, human foreskin fibroblasts, one type of somatic cell, use non-homologous end joining (NHEJ), an error prone DNA repair process, as the primary pathway for repairing double-strand breaks (DSBs) during all cell cycle stages. Because of its error-prone nature, NHEJ tends to produce mutations in a cell's clonal descendants.
ES cells use a different strategy to deal with DSBs. Because ES cells give rise to all of the cell types of an organism including the cells of the germ line, mutations arising in ES cells due to faulty DNA repair are a more serious problem than in differentiated somatic cells. Consequently, robust mechanisms are needed in ES cells to repair DNA damages accurately, and if repair fails, to remove those cells with un-repaired DNA damages. Thus, mouse ES cells predominantly use high fidelity homologous recombinational repair (HRR) to repair DSBs. This type of repair depends on the interaction of the two sister chromosomes formed during S phase and present together during the G2 phase of the cell cycle. HRR can accurately repair DSBs in one sister chromosome by using intact information from the other sister chromosome. Cells in the G1 phase of the cell cycle (i.e. after metaphase/cell division but prior the next round of replication) have only one copy of each chromosome (i.e. sister chromosomes aren't present). Mouse ES cells lack a G1 checkpoint and do not undergo cell cycle arrest upon acquiring DNA damage. Rather they undergo programmed cell death (apoptosis) in response to DNA damage. Apoptosis can be used as a fail-safe strategy to remove cells with un-repaired DNA damages in order to avoid mutation and progression to cancer. Consistent with this strategy, mouse ES stem cells have a mutation frequency about 100-fold lower than that of isogenic mouse somatic cells.
Clinical trial
On January 23, 2009, Phase I clinical trials for transplantation of oligodendrocytes (a cell type of the brain and spinal cord) derived from human ESCs into spinal cord-injured individuals received approval from the U.S. Food and Drug Administration (FDA), marking it the world's first human ESC human trial. The study leading to this scientific advancement was conducted by Hans Keirstead and colleagues at the University of California, Irvine and supported by Geron Corporation of Menlo Park, CA, founded by Michael D. West, PhD. A previous experiment had shown an improvement in locomotor recovery in spinal cord-injured rats after a 7-day delayed transplantation of human ESCs that had been pushed into an oligodendrocytic lineage. The phase I clinical study was designed to enroll about eight to ten paraplegics who have had their injuries no longer than two weeks before the trial begins, since the cells must be injected before scar tissue is able to form. The researchers emphasized that the injections were not expected to fully cure the patients and restore all mobility. Based on the results of the rodent trials, researchers speculated that restoration of myelin sheathes and an increase in mobility might occur. This first trial was primarily designed to test the safety of these procedures and if everything went well, it was hoped that it would lead to future studies that involve people with more severe disabilities. The trial was put on hold in August 2009 due to FDA concerns regarding a small number of microscopic cysts found in several treated rat models but the hold was lifted on July 30, 2010.
In October 2010 researchers enrolled and administered ESCs to the first patient at Shepherd Center in Atlanta. The makers of the stem cell therapy, Geron Corporation, estimated that it would take several months for the stem cells to replicate and for the GRNOPC1 therapy to be evaluated for success or failure.
In November 2011 Geron announced it was halting the trial and dropping out of stem cell research for financial reasons, but would continue to monitor existing patients, and was attempting to find a partner that could continue their research. In 2013 BioTime, led by CEO Dr. Michael D. West, acquired all of Geron's stem cell assets, with the stated intention of restarting Geron's embryonic stem cell-based clinical trial for spinal cord injury research.
BioTime company Asterias Biotherapeutics (NYSE MKT: AST) was granted a $14.3 million Strategic Partnership Award by the California Institute for Regenerative Medicine (CIRM) to re-initiate the world's first embryonic stem cell-based human clinical trial, for spinal cord injury. Supported by California public funds, CIRM is the largest funder of stem cell-related research and development in the world.
The award provides funding for Asterias to reinitiate clinical development of AST-OPC1 in subjects with spinal cord injury and to expand clinical testing of escalating doses in the target population intended for future pivotal trials.
AST-OPC1 is a population of cells derived from human embryonic stem cells (hESCs) that contains oligodendrocyte progenitor cells (OPCs). OPCs and their mature derivatives called oligodendrocytes provide critical functional support for nerve cells in the spinal cord and brain. Asterias recently presented the results from phase 1 clinical trial testing of a low dose of AST-OPC1 in patients with neurologically complete thoracic spinal cord injury. The results showed that AST-OPC1 was successfully delivered to the injured spinal cord site. Patients followed 2–3 years after AST-OPC1 administration showed no evidence of serious adverse events associated with the cells in detailed follow-up assessments including frequent neurological exams and MRIs. Immune monitoring of subjects through one year post-transplantation showed no evidence of antibody-based or cellular immune responses to AST-OPC1. In four of the five subjects, serial MRI scans performed throughout the 2–3 year follow-up period indicate that reduced spinal cord cavitation may have occurred and that AST-OPC1 may have had some positive effects in reducing spinal cord tissue deterioration. There was no unexpected neurological degeneration or improvement in the five subjects in the trial as evaluated by the International Standards for Neurological Classification of Spinal Cord Injury (ISNCSCI) exam.
The Strategic Partnership III grant from CIRM will provide funding to Asterias to support the next clinical trial of AST-OPC1 in subjects with spinal cord injury, and for Asterias' product development efforts to refine and scale manufacturing methods to support later-stage trials and eventually commercialization. CIRM funding will be conditional on FDA approval for the trial, completion of a definitive agreement between Asterias and CIRM, and Asterias' continued progress toward the achievement of certain pre-defined project milestones.
Concern and controversy
Adverse effects
The major concern with the possible transplantation of ESCs into patients as therapies is their ability to form tumors including teratomas. Safety issues prompted the FDA to place a hold on the first ESC clinical trial, however no tumors were observed.
The main strategy to enhance the safety of ESCs for potential clinical use is to differentiate the ESCs into specific cell types (e.g. neurons, muscle, liver cells) that have reduced or eliminated ability to cause tumors. Following differentiation, the cells are subjected to sorting by flow cytometry for further purification. ESCs are predicted to be inherently safer than iPS cells created with genetically integrating viral vectors because they are not genetically modified with genes such as c-Myc that are linked to cancer. Nonetheless, ESCs express very high levels of the iPS inducing genes and these genes including Myc are essential for ESC self-renewal and pluripotency, and potential strategies to improve safety by eliminating c-Myc expression are unlikely to preserve the cells' "stemness". However, N-myc and L-myc have been identified to induce iPS cells instead of c-myc with similar efficiency. Later protocols to induce pluripotency bypass these problems completely by using non-integrating RNA viral vectors such as sendai virus or mRNA transfection.
Ethical debate
Due to the nature of embryonic stem cell research, there are a lot of controversial opinions on the topic. Since harvesting embryonic stem cells usually necessitates destroying the embryo from which those cells are obtained, the moral status of the embryo comes into question. Some people claim that the embryo is too young to achieve personhood or that the embryo, if donated from an IVF clinic (where labs typically acquire embryos), would otherwise go to medical waste anyway. Opponents of ESC research claim that an embryo is a human life, therefore destroying it is murder and the embryo must be protected under the same ethical view as a more developed human being.
History
1964: Lewis Kleinsmith and G. Barry Pierce Jr. isolated a single type of cell from a teratocarcinoma, a tumor now known from a germ cell. These cells were isolated from the teratocarcinoma replicated and grew in cell culture as a stem cell and are now known as embryonal carcinoma (EC) cells. Although similarities in morphology and differentiating potential (pluripotency) led to the use of EC cells as the in vitro model for early mouse development, EC cells harbor genetic mutations and often abnormal karyotypes that accumulated during the development of the teratocarcinoma. These genetic aberrations further emphasized the need to be able to culture pluripotent cells directly from the inner cell mass.
1981: Embryonic stem cells (ES cells) were independently first derived from a mouse embryos by two groups. Martin Evans and Matthew Kaufman from the Department of Genetics, University of Cambridge published first in July, revealing a new technique for culturing the mouse embryos in the uterus to allow for an increase in cell number, allowing for the derivation of ES cell from these embryos. Gail R. Martin, from the Department of Anatomy, University of California, San Francisco, published her paper in December and coined the term "Embryonic Stem Cell". She showed that embryos could be cultured in vitro and that ES cells could be derived from these embryos.
1989: Mario R. Cappechi, Martin J. Evans, and Oliver Smithies publish their research that details their isolation and genetic modifications of embryonic stem cells, creating the first "knockout mice". In creating knockout mice, this publication provided scientists with an entirely new way to study disease.
1996: Dolly, was the first mammal cloned from an adult cell by the Roslin Institute of the University of Edinburgh. This experiment instituted the proposition that specialized adult cells obtain the genetic makeup to perform a specific task; which established a basis for further research within a variety of cloning techniques. The Dolly experiment was performed by obtaining the mammalian udder cells from a sheep (Dolly) and differentiating these cells until division was concluded. An egg cell was then procured from a different sheep host and the nucleus was removed. An udder cell was placed next to the egg cell and connected by electricity causing this cell to share DNA. This egg cell differentiated into an embryo and the embryo was inserted into a third sheep which gave birth to the clone version of Dolly.
1998: A team from the University of Wisconsin, Madison (James A. Thomson, Joseph Itskovitz-Eldor, Sander S. Shapiro, Michelle A. Waknitz, Jennifer J. Swiergiel, Vivienne S. Marshall, and Jeffrey M. Jones) publish a paper titled "Embryonic Stem Cell Lines Derived From Human Blastocysts". The researchers behind this study not only created the first embryonic stem cells, but recognized their pluripotency, as well as their capacity for self-renewal. The abstract of the paper notes the significance of the discovery with regards to the fields of developmental biology and drug discovery.
2001: President George W. Bush allows federal funding to support research on roughly 60—at this time, already existing—lines of embryonic stem cells. Seeing as the limited lines that Bush allowed research on had already been established, this law supported embryonic stem cell research without raising any ethical questions that could arise with the creation of new lines under federal budget.
2006: Japanese scientists Shinya Yamanaka and Kazutoshi Takashi publish a paper describing the induction of pluripotent stem cells from cultures of adult mouse fibroblasts. Induced pluripotent stem cells (iPSCs) are a huge discovery, as they are seemingly identical to embryonic stem cells and could be used without sparking the same moral controversy.
January, 2009: The US Food and Drug Administration (FDA) provides approval for Geron Corporation's phase I trial of their human embryonic stem cell-derived treatment for spinal cord injuries. The announcement was met with excitement from the scientific community, but also with wariness from stem cell opposers. The treatment cells were, however, derived from the cell lines approved under George W. Bush's ESC policy.
March, 2009: Executive Order 13505 is signed by President Barack Obama, removing the restrictions put in place on federal funding for human stem cells by the previous presidential administration. This would allow the National Institutes of Health (NIH) to provide funding for hESC research. The document also states that the NIH must provide revised federal funding guidelines within 120 days of the order's signing.
Techniques and conditions for derivation and culture
Derivation from humans
In vitro fertilization generates multiple embryos. The surplus of embryos is not clinically used or is unsuitable for implantation into the patient, and therefore may be donated by the donor with consent. Human embryonic stem cells can be derived from these donated embryos or additionally they can also be extracted from cloned embryos created using a cell from a patient and a donated egg through the process of somatic cell nuclear transfer. The inner cell mass (cells of interest), from the blastocyst stage of the embryo, is separated from the trophectoderm, the cells that would differentiate into extra-embryonic tissue. Immunosurgery, the process in which antibodies are bound to the trophectoderm and removed by another solution, and mechanical dissection are performed to achieve separation. The resulting inner cell mass cells are plated onto cells that will supply support. The inner cell mass cells attach and expand further to form a human embryonic cell line, which are undifferentiated. These cells are fed daily and are enzymatically or mechanically separated every four to seven days. For differentiation to occur, the human embryonic stem cell line is removed from the supporting cells to form embryoid bodies, is co-cultured with a serum containing necessary signals, or is grafted in a three-dimensional scaffold to result.
Derivation from other animals
Embryonic stem cells are derived from the inner cell mass of the early embryo, which are harvested from the donor mother animal. Martin Evans and Matthew Kaufman reported a technique that delays embryo implantation, allowing the inner cell mass to increase. This process includes removing the donor mother's ovaries and dosing her with progesterone, changing the hormone environment, which causes the embryos to remain free in the uterus. After 4–6 days of this intrauterine culture, the embryos are harvested and grown in in vitro culture until the inner cell mass forms “egg cylinder-like structures,” which are dissociated into single cells, and plated on fibroblasts treated with mitomycin-c (to prevent fibroblast mitosis). Clonal cell lines are created by growing up a single cell. Evans and Kaufman showed that the cells grown out from these cultures could form teratomas and embryoid bodies, and differentiate in vitro, all of which indicating that the cells are pluripotent.
Gail Martin derived and cultured her ES cells differently. She removed the embryos from the donor mother at approximately 76 hours after copulation and cultured them overnight in a medium containing serum. The following day, she removed the inner cell mass from the late blastocyst using microsurgery. The extracted inner cell mass was cultured on fibroblasts treated with mitomycin-c in a medium containing serum and conditioned by ES cells. After approximately one week, colonies of cells grew out. These cells grew in culture and demonstrated pluripotent characteristics, as demonstrated by the ability to form teratomas, differentiate in vitro, and form embryoid bodies. Martin referred to these cells as ES cells.
It is now known that the feeder cells provide leukemia inhibitory factor (LIF) and serum provides bone morphogenetic proteins (BMPs) that are necessary to prevent ES cells from differentiating. These factors are extremely important for the efficiency of deriving ES cells. Furthermore, it has been demonstrated that different mouse strains have different efficiencies for isolating ES cells. Current uses for mouse ES cells include the generation of transgenic mice, including knockout mice. For human treatment, there is a need for patient specific pluripotent cells. Generation of human ES cells is more difficult and faces ethical issues. So, in addition to human ES cell research, many groups are focused on the generation of induced pluripotent stem cells (iPS cells).
Potential methods for new cell line derivation
On August 23, 2006, the online edition of Nature scientific journal published a letter by Dr. Robert Lanza (medical director of Advanced Cell Technology in Worcester, MA) stating that his team had found a way to extract embryonic stem cells without destroying the actual embryo. This technical achievement would potentially enable scientists to work with new lines of embryonic stem cells derived using public funding in the US, where federal funding was at the time limited to research using embryonic stem cell lines derived prior to August 2001. In March, 2009, the limitation was lifted.
Human embryonic stem cells have also been derived by somatic cell nuclear transfer (SCNT). This approach has also sometimes been referred to as "therapeutic cloning" because SCNT bears similarity to other kinds of cloning in that nuclei are transferred from a somatic cell into an enucleated zygote. However, in this case SCNT was used to produce embryonic stem cell lines in a lab, not living organisms via a pregnancy. The "therapeutic" part of the name is included because of the hope that SCNT produced embryonic stem cells could have clinical utility.
Induced pluripotent stem cells
The iPS cell technology was pioneered by Shinya Yamanaka's lab in Kyoto, Japan, who showed in 2006 that the introduction of four specific genes encoding transcription factors could convert adult cells into pluripotent stem cells. He was awarded the 2012 Nobel Prize along with Sir John Gurdon "for the discovery that mature cells can be reprogrammed to become pluripotent."
In 2007, it was shown that pluripotent stem cells, highly similar to embryonic stem cells, can be induced by the delivery of four factors (Oct3/4, Sox2, c-Myc, and Klf4) to differentiated cells. Utilizing the four genes previously listed, the differentiated cells are "reprogrammed" into pluripotent stem cells, allowing for the generation of pluripotent/embryonic stem cells without the embryo. The morphology and growth factors of these lab induced pluripotent cells, are equivalent to embryonic stem cells, leading these cells to be known as induced pluripotent stem cells (iPS cells). This observation was observed in mouse pluripotent stem cells, originally, but now can be performed in human adult fibroblasts using the same four genes.
Because ethical concerns regarding embryonic stem cells typically are about their derivation from terminated embryos, it is believed that reprogramming to these iPS cells may be less controversial.
This may enable the generation of patient specific ES cell lines that could potentially be used for cell replacement therapies. In addition, this will allow the generation of ES cell lines from patients with a variety of genetic diseases and will provide invaluable models to study those diseases.
However, as a first indication that the iPS cell technology can in rapid succession lead to new cures, it was used by a research team headed by Rudolf Jaenisch of the Whitehead Institute for Biomedical Research in Cambridge, Massachusetts, to cure mice of sickle cell anemia, as reported by Science journal's online edition on December 6, 2007.
On January 16, 2008, a California-based company, Stemagen, announced that they had created the first mature cloned human embryos from single skin cells taken from adults. These embryos can be harvested for patient matching embryonic stem cells.
Contamination by reagents used in cell culture
The online edition of Nature Medicine published a study on January 24, 2005, which stated that the human embryonic stem cells available for federally funded research are contaminated with non-human molecules from the culture medium used to grow the cells. It is a common technique to use mouse cells and other animal cells to maintain the pluripotency of actively dividing stem cells. The problem was discovered when non-human sialic acid in the growth medium was found to compromise the potential uses of the embryonic stem cells in humans, according to scientists at the University of California, San Diego.
However, a study published in the online edition of Lancet Medical Journal on March 8, 2005, detailed information about a new stem cell line that was derived from human embryos under completely cell- and serum-free conditions. After more than 6 months of undifferentiated proliferation, these cells demonstrated the potential to form derivatives of all three embryonic germ layers both in vitro and in teratomas. These properties were also successfully maintained (for more than 30 passages) with the established stem cell lines.
Muse cells
Muse cells (Multi-lineage differentiating stress enduring cell) are non-cancerous pluripotent stem cell found in adults. They were discovered in 2010 by Mari Dezawa and her research group. Muse cells reside in the connective tissue of nearly every organ including the umbilical cord, bone marrow and peripheral blood. They are collectable from commercially obtainable mesenchymal cells such as human fibroblasts, bone marrow-mesenchymal stem cells and adipose-derived stem cells. Muse cells are able to generate cells representative of all three germ layers from a single cell both spontaneously and under cytokine induction. Expression of pluripotency genes and triploblastic differentiation are self-renewable over generations. Muse cells do not undergo teratoma formation when transplanted into a host environment in vivo, eradicating the risk of tumorigenesis through unbridled cell proliferation.
| Biology and health sciences | Cell processes | Biology |
1029211 | https://en.wikipedia.org/wiki/Metabolomics | Metabolomics | Metabolomics is the scientific study of chemical processes involving metabolites, the small molecule substrates, intermediates, and products of cell metabolism. Specifically, metabolomics is the "systematic study of the unique chemical fingerprints that specific cellular processes leave behind", the study of their small-molecule metabolite profiles. The metabolome represents the complete set of metabolites in a biological cell, tissue, organ, or organism, which are the end products of cellular processes. Messenger RNA (mRNA), gene expression data, and proteomic analyses reveal the set of gene products being produced in the cell, data that represents one aspect of cellular function. Conversely, metabolic profiling can give an instantaneous snapshot of the physiology of that cell, and thus, metabolomics provides a direct "functional readout of the physiological state" of an organism. There are indeed quantifiable correlations between the metabolome and the other cellular ensembles (genome, transcriptome, proteome, and lipidome), which can be used to predict metabolite abundances in biological samples from, for example mRNA abundances. One of the ultimate challenges of systems biology is to integrate metabolomics with all other -omics information to provide a better understanding of cellular biology.
History
The concept that individuals might have a "metabolic profile" that could be reflected in the makeup of their biological fluids was introduced by Roger Williams in the late 1940s, who used paper chromatography to suggest characteristic metabolic patterns in urine and saliva were associated with diseases such as schizophrenia. However, it was only through technological advancements in the 1960s and 1970s that it became feasible to quantitatively (as opposed to qualitatively) measure metabolic profiles. The term "metabolic profile" was introduced by Horning, et al. in 1971 after they demonstrated that gas chromatography-mass spectrometry (GC-MS) could be used to measure compounds present in human urine and tissue extracts. The Horning group, along with that of Linus Pauling and Arthur B. Robinson led the development of GC-MS methods to monitor the metabolites present in urine through the 1970s.
Concurrently, NMR spectroscopy, which was discovered in the 1940s, was also undergoing rapid advances. In 1974, Seeley et al. demonstrated the utility of using NMR to detect metabolites in unmodified biological samples. This first study on muscle highlighted the value of NMR in that it was determined that 90% of cellular ATP is complexed with magnesium. As sensitivity has improved with the evolution of higher magnetic field strengths and magic angle spinning, NMR continues to be a leading analytical tool to investigate metabolism. Recent efforts to utilize NMR for metabolomics have been largely driven by the laboratory of Jeremy K. Nicholson at Birkbeck College, University of London and later at Imperial College London. In 1984, Nicholson showed 1H NMR spectroscopy could potentially be used to diagnose diabetes mellitus, and later pioneered the application of pattern recognition methods to NMR spectroscopic data.
In 1994 and 1996, liquid chromatography mass spectrometry metabolomics experiments were performed by Gary Siuzdak while working with Richard Lerner (then president of the Scripps Research Institute) and Benjamin Cravatt, to analyze the cerebral spinal fluid from sleep deprived animals. One molecule of particular interest, oleamide, was observed and later shown to have sleep inducing properties. This work is one of the earliest such experiments combining liquid chromatography and mass spectrometry in metabolomics.
In 2005, the first metabolomics tandem mass spectrometry database, METLIN, for characterizing human metabolites was developed in the Siuzdak laboratory at the Scripps Research Institute. METLIN has since grown and as of December, 2023, METLIN contains MS/MS experimental data on over 930,000 molecular standards and other chemical entities, each compound having experimental tandem mass spectrometry data generated from molecular standards at multiple collision energies and in positive and negative ionization modes. METLIN is the largest repository of tandem mass spectrometry data of its kind. The dedicated academic journal Metabolomics first appeared in 2005, founded by its current editor-in-chief Roy Goodacre.
In 2005, the Siuzdak lab was engaged in identifying metabolites associated with sepsis and in an effort to address the issue of statistically identifying the most relevant dysregulated metabolites across hundreds of LC/MS datasets, the first algorithm was developed to allow for the nonlinear alignment of mass spectrometry metabolomics data. Called XCMS, it has since (2012) been developed as an online tool and as of 2019 (with METLIN) has over 30,000 registered users.
On 23 January 2007, the Human Metabolome Project, led by David S. Wishart, completed the first draft of the human metabolome, consisting of a database of approximately 2,500 metabolites, 1,200 drugs and 3,500 food components. Similar projects have been underway in several plant species, most notably Medicago truncatula and Arabidopsis thaliana for several years.
As late as mid-2010, metabolomics was still considered an "emerging field". Further, it was noted that further progress in the field depended in large part, through addressing otherwise "irresolvable technical challenges", by technical evolution of mass spectrometry instrumentation.
In 2015, real-time metabolome profiling was demonstrated for the first time.
Metabolome
The metabolome refers to the complete set of small-molecule (<1.5 kDa) metabolites (such as metabolic intermediates, hormones and other signaling molecules, and secondary metabolites) to be found within a biological sample, such as a single organism. The word was coined in analogy with transcriptomics and proteomics; like the transcriptome and the proteome, the metabolome is dynamic, changing from second to second. Although the metabolome can be defined readily enough, it is not currently possible to analyse the entire range of metabolites by a single analytical method.
In January 2007, scientists at the University of Alberta and the University of Calgary completed the first draft of the human metabolome. The Human Metabolome Database (HMDB) is perhaps the most extensive public metabolomic spectral database to date and is a freely available electronic database (www.hmdb.ca) containing detailed information about small molecule metabolites found in the human body. It is intended to be used for applications in metabolomics, clinical chemistry, biomarker discovery and general education. The database is designed to contain or link three kinds of data:
Chemical data,
Clinical data and
Molecular biology/biochemistry data.
The database contains 220,945 metabolite entries including both water-soluble and lipid soluble metabolites. Additionally, 8,610 protein sequences (enzymes and transporters) are linked to these metabolite entries. Each MetaboCard entry contains 130 data fields with 2/3 of the information being devoted to chemical/clinical data and the other 1/3 devoted to enzymatic or biochemical data. The version 3.5 of the HMDB contains >16,000 endogenous metabolites, >1,500 drugs and >22,000 food constituents or food metabolites. This information, available at the Human Metabolome Database and based on analysis of information available in the current scientific literature, is far from complete. In contrast, much more is known about the metabolomes of other organisms. For example, over 50,000 metabolites have been characterized from the plant kingdom, and many thousands of metabolites have been identified and/or characterized from single plants.
Each type of cell and tissue has a unique metabolic ‘fingerprint’ that can elucidate organ or tissue-specific information. Bio-specimens used for metabolomics analysis include but not limit to plasma, serum, urine, saliva, feces, muscle, sweat, exhaled breath and gastrointestinal fluid. The ease of collection facilitates high temporal resolution, and because they are always at dynamic equilibrium with the body, they can describe the host as a whole. Genome can tell what could happen, transcriptome can tell what appears to be happening, proteome can tell what makes it happen and metabolome can tell what has happened and what is happening.
Metabolites
Metabolites are the substrates, intermediates and products of metabolism. Within the context of metabolomics, a metabolite is usually defined as any molecule less than 1.5 kDa in size. However, there are exceptions to this depending on the sample and detection method. For example, macromolecules such as lipoproteins and albumin are reliably detected in NMR-based metabolomics studies of blood plasma. In plant-based metabolomics, it is common to refer to "primary" and "secondary" metabolites. A primary metabolite is directly involved in the normal growth, development, and reproduction. A secondary metabolite is not directly involved in those processes, but usually has important ecological function. Examples include antibiotics and pigments. By contrast, in human-based metabolomics, it is more common to describe metabolites as being either endogenous (produced by the host organism) or exogenous. Metabolites of foreign substances such as drugs are termed xenometabolites.
The metabolome derives from a large network of metabolic reactions, where outputs from one enzymatic chemical reaction are inputs to other chemical reactions. Such systems have been described as hypercycles.
Metabonomics
Metabonomics is defined as "the quantitative measurement of the dynamic multiparametric metabolic response of living systems to pathophysiological stimuli or genetic modification". The word origin is from the Greek μεταβολή meaning change and nomos meaning a rule set or set of laws. This approach was pioneered by Jeremy Nicholson at Murdoch University and has been used in toxicology, disease diagnosis and a number of other fields. Historically, the metabonomics approach was one of the first methods to apply the scope of systems biology to studies of metabolism.
There has been some disagreement over the exact differences between 'metabolomics' and 'metabonomics'. The difference between the two terms is not related to choice of analytical platform: although metabonomics is more associated with NMR spectroscopy and metabolomics with mass spectrometry-based techniques, this is simply because of usages amongst different groups that have popularized the different terms. While there is still no absolute agreement, there is a growing consensus that 'metabolomics' places a greater emphasis on metabolic profiling at a cellular or organ level and is primarily concerned with normal endogenous metabolism. 'Metabonomics' extends metabolic profiling to include information about perturbations of metabolism caused by environmental factors (including diet and toxins), disease processes, and the involvement of extragenomic influences, such as gut microflora. This is not a trivial difference; metabolomic studies should, by definition, exclude metabolic contributions from extragenomic sources, because these are external to the system being studied. However, in practice, within the field of human disease research there is still a large degree of overlap in the way both terms are used, and they are often in effect synonymous.
Exometabolomics
Exometabolomics, or "metabolic footprinting", is the study of extracellular metabolites. It uses many techniques from other subfields of metabolomics, and has applications in biofuel development, bioprocessing, determining drugs' mechanism of action, and studying intercellular interactions.
Analytical technologies
The typical workflow of metabolomics studies is shown in the figure. First, samples are collected from tissue, plasma, urine, saliva, cells, etc. Next, metabolites extracted often with the addition of internal standards and derivatization. During sample analysis, metabolites are quantified (liquid chromatography or gas chromatography coupled with MS and/or NMR spectroscopy). The raw output data can be used for metabolite feature extraction and further processed before statistical analysis (such as principal component analysis, PCA). Many bioinformatic tools and software are available to identify associations with disease states and outcomes, determine significant correlations, and characterize metabolic signatures with existing biological knowledge.
Separation methods
Initially, analytes in a metabolomic sample comprise a highly complex mixture. This complex mixture can be simplified prior to detection by separating some analytes from others. Separation achieves various goals: analytes which cannot be resolved by the detector may be separated in this step; in MS analysis, ion suppression is reduced; the retention time of the analyte serves as information regarding its identity. This separation step is not mandatory and is often omitted in NMR and "shotgun" based approaches such as shotgun lipidomics.
Gas chromatography (GC), especially when interfaced with mass spectrometry (GC-MS), is a widely used separation technique for metabolomic analysis. GC offers very high chromatographic resolution, and can be used in conjunction with a flame ionization detector (GC/FID) or a mass spectrometer (GC-MS). The method is especially useful for identification and quantification of small and volatile molecules. However, a practical limitation of GC is the requirement of chemical derivatization for many biomolecules as only volatile chemicals can be analysed without derivatization. In cases where greater resolving power is required, two-dimensional chromatography (GCxGC) can be applied.
High performance liquid chromatography (HPLC) has emerged as the most common separation technique for metabolomic analysis. With the advent of electrospray ionization, HPLC was coupled to MS. In contrast with GC, HPLC has lower chromatographic resolution, but requires no derivatization for polar molecules, and separates molecules in the liquid phase. Additionally HPLC has the advantage that a much wider range of analytes can be measured with a higher sensitivity than GC methods.
Capillary electrophoresis (CE) has a higher theoretical separation efficiency than HPLC (although requiring much more time per separation), and is suitable for use with a wider range of metabolite classes than is GC. As for all electrophoretic techniques, it is most appropriate for charged analytes.
In direct-infusion mass spectrometry (DI-MS), sample is directly introduced into the spectrometer and separation steps are skipped. DI-MS can be employed to perform single cell metabolic analysis of human cells.
Detection methods
Mass spectrometry (MS) is used to identify and quantify metabolites after optional separation by GC, HPLC, or CE. GC-MS was the first hyphenated technique to be developed. Identification leverages the distinct patterns in which analytes fragment. These patterns can be thought of as a mass spectral fingerprint. Libraries exist that allow identification of a metabolite according to this fragmentation pattern . MS is both sensitive and can be very specific. There are also a number of techniques which use MS as a stand-alone technology: the sample is infused directly into the mass spectrometer with no prior separation, and the MS provides sufficient selectivity to both separate and to detect metabolites.
For analysis by mass spectrometry, the analytes must be imparted with a charge and transferred to the gas phase. Electron ionization (EI) is the most common ionization technique applied to GC separations as it is amenable to low pressures. EI also produces fragmentation of the analyte, both providing structural information while increasing the complexity of the data and possibly obscuring the molecular ion. Atmospheric-pressure chemical ionization (APCI) is an atmospheric pressure technique that can be applied to all the above separation techniques. APCI is a gas phase ionization method, which provides slightly more aggressive ionization than ESI which is suitable for less polar compounds. Electrospray ionization (ESI) is the most common ionization technique applied in LC/MS. This soft ionization is most successful for polar molecules with ionizable functional groups. Another commonly used soft ionization technique is secondary electrospray ionization (SESI).
In the 2000s, surface-based mass analysis has seen a resurgence, with new MS technologies focused on increasing sensitivity, minimizing background, and reducing sample preparation. The ability to analyze metabolites directly from biofluids and tissues continues to challenge current MS technology, largely because of the limits imposed by the complexity of these samples, which contain thousands to tens of thousands of metabolites. Among the technologies being developed to address this challenge is Nanostructure-Initiator MS (NIMS), a desorption/ ionization approach that does not require the application of matrix and thereby facilitates small-molecule (i.e., metabolite) identification. MALDI is also used; however, the application of a MALDI matrix can add significant background at that complicates analysis of the low-mass range (i.e., metabolites). In addition, the size of the resulting matrix crystals limits the spatial resolution that can be achieved in tissue imaging. Because of these limitations, several other matrix-free desorption/ionization approaches have been applied to the analysis of biofluids and tissues.
Secondary ion mass spectrometry (SIMS) was one of the first matrix-free desorption/ionization approaches used to analyze metabolites from biological samples. SIMS uses a high-energy primary ion beam to desorb and generate secondary ions from a surface. The primary advantage of SIMS is its high spatial resolution (as small as 50 nm), a powerful characteristic for tissue imaging with MS. However, SIMS has yet to be readily applied to the analysis of biofluids and tissues because of its limited sensitivity at and analyte fragmentation generated by the high-energy primary ion beam. Desorption electrospray ionization (DESI) is a matrix-free technique for analyzing biological samples that uses a charged solvent spray to desorb ions from a surface. Advantages of DESI are that no special surface is required and the analysis is performed at ambient pressure with full access to the sample during acquisition. A limitation of DESI is spatial resolution because "focusing" the charged solvent spray is difficult. However, a recent development termed laser ablation ESI (LAESI) is a promising approach to circumvent this limitation. Most recently, ion trap techniques such as orbitrap mass spectrometry are also applied to metabolomics research.
Nuclear magnetic resonance (NMR) spectroscopy is the only detection technique which does not rely on separation of the analytes, and the sample can thus be recovered for further analyses. All kinds of small molecule metabolites can be measured simultaneously - in this sense, NMR is close to being a universal detector. The main advantages of NMR are high analytical reproducibility and simplicity of sample preparation. Practically, however, it is relatively insensitive compared to mass spectrometry-based techniques.
Although NMR and MS are the most widely used modern-day techniques for detection, there are other methods in use. These include Fourier-transform ion cyclotron resonance, ion-mobility spectrometry, electrochemical detection (coupled to HPLC), Raman spectroscopy and radiolabel (when combined with thin-layer chromatography).
Statistical methods
The data generated in metabolomics usually consist of measurements performed on subjects under various conditions. These measurements may be digitized spectra, or a list of metabolite features. In its simplest form, this generates a matrix with rows corresponding to subjects and columns corresponding with metabolite features (or vice versa). Several statistical programs are currently available for analysis of both NMR and mass spectrometry data. A great number of free software are already available for the analysis of metabolomics data shown in the table. Some statistical tools listed in the table were designed for NMR data analyses were also useful for MS data. For mass spectrometry data, software is available that identifies molecules that vary in subject groups on the basis of mass-over-charge value and sometimes retention time depending on the experimental design.
Once metabolite data matrix is determined, unsupervised data reduction techniques (e.g. PCA) can be used to elucidate patterns and connections. In many studies, including those evaluating drug-toxicity and some disease models, the metabolites of interest are not known a priori. This makes unsupervised methods, those with no prior assumptions of class membership, a popular first choice. The most common of these methods includes principal component analysis (PCA) which can efficiently reduce the dimensions of a dataset to a few which explain the greatest variation. When analyzed in the lower-dimensional PCA space, clustering of samples with similar metabolic fingerprints can be detected. PCA algorithms aim to replace all correlated variables with a much smaller number of uncorrelated variables (referred to as principal components (PCs)) and retain most of the information in the original dataset. This clustering can elucidate patterns and assist in the determination of disease biomarkers – metabolites that correlate most with class membership.
Linear models are commonly used for metabolomics data, but are affected by multicollinearity. On the other hand, multivariate statistics are thriving methods for high-dimensional correlated metabolomics data, of which the most popular one is Projection to Latent Structures (PLS) regression and its classification version PLS-DA. Other data mining methods, such as random forest, support-vector machines, etc. are received increasing attention for untargeted metabolomics data analysis. In the case of univariate methods, variables are analyzed one by one using classical statistics tools (such as Student's t-test, ANOVA or mixed models) and only these with sufficient small p-values are considered relevant. However, correction strategies should be used to reduce false discoveries when multiple comparisons are conducted since there is no standard method for measuring the total amount of metabolites directly in untargeted metabolomics. For multivariate analysis, models should always be validated to ensure that the results can be generalized.
Machine learning and data mining
Machine learning is a powerful tool that can be used in metabolomics analysis. Recently, scientists have developed retention time prediction software. These tools allow researchers to apply artificial intelligence to the retention time prediction of small molecules in complex mixture, such as human plasma, plant extracts, foods, or microbial cultures. Retention time prediction increases the identification rate in liquid chromatography and can lead to an improved biological interpretation of metabolomics data.
Key applications
Toxicity assessment/toxicology by metabolic profiling (especially of urine or blood plasma samples) detects the physiological changes caused by toxic insult of a chemical (or mixture of chemicals). In many cases, the observed changes can be related to specific syndromes, e.g. a specific lesion in liver or kidney. This is of particular relevance to pharmaceutical companies wanting to test the toxicity of potential drug candidates: if a compound can be eliminated before it reaches clinical trials on the grounds of adverse toxicity, it saves the enormous expense of the trials.
For functional genomics, metabolomics can be an excellent tool for determining the phenotype caused by a genetic manipulation, such as gene deletion or insertion. Sometimes this can be a sufficient goal in itself—for instance, to detect any phenotypic changes in a genetically modified plant intended for human or animal consumption. More exciting is the prospect of predicting the function of unknown genes by comparison with the metabolic perturbations caused by deletion/insertion of known genes. Such advances are most likely to come from model organisms such as Saccharomyces cerevisiae and Arabidopsis thaliana. The Cravatt laboratory at the Scripps Research Institute has recently applied this technology to mammalian systems, identifying the N-acyltaurines as previously uncharacterized endogenous substrates for the enzyme fatty acid amide hydrolase (FAAH) and the monoalkylglycerol ethers (MAGEs) as endogenous substrates for the uncharacterized hydrolase KIAA1363.
Metabologenomics is a novel approach to integrate metabolomics and genomics data by correlating microbial-exported metabolites with predicted biosynthetic genes. This bioinformatics-based pairing method enables natural product discovery at a larger-scale by refining non-targeted metabolomic analyses to identify small molecules with related biosynthesis and to focus on those that may not have previously well known structures.
Fluxomics is a further development of metabolomics. The disadvantage of metabolomics is that it only provides the user with abundances or concentrations of metabolites, while fluxomics determines the reaction rates of metabolic reactions and can trace metabolites in a biological system over time.
Nutrigenomics is a generalised term which links genomics, transcriptomics, proteomics and metabolomics to human nutrition. In general, in a given body fluid, a metabolome is influenced by endogenous factors such as age, sex, body composition and genetics as well as underlying pathologies. The large bowel microflora are also a very significant potential confounder of metabolic profiles and could be classified as either an endogenous or exogenous factor. The main exogenous factors are diet and drugs. Diet can then be broken down to nutrients and non-nutrients. Metabolomics is one means to determine a biological endpoint, or metabolic fingerprint, which reflects the balance of all these forces on an individual's metabolism. Thanks to recent cost reductions, metabolomics has now become accessible for companion animals, such as pregnant dogs.
Plant metabolomics is designed to study the overall changes in metabolites of plant samples and then conduct deep data mining and chemometric analysis. Specialized metabolites are considered components of plant defense systems biosynthesized in response to biotic and abiotic stresses. Metabolomics approaches have recently been used to assess the natural variance in metabolite content between individual plants, an approach with great potential for the improvement of the compositional quality of crops.
| Biology and health sciences | Metabolic processes | Biology |
1029967 | https://en.wikipedia.org/wiki/Hypernucleus | Hypernucleus | A hypernucleus is similar to a conventional atomic nucleus, but contains at least one hyperon in addition to the normal protons and neutrons. Hyperons are a category of baryon particles that carry non-zero strangeness quantum number, which is conserved by the strong and electromagnetic interactions.
A variety of reactions give access to depositing one or more units of strangeness in a nucleus. Hypernuclei containing the lightest hyperon, the lambda (Λ), tend to be more tightly bound than normal nuclei, though they can decay via the weak force with a mean lifetime of around . Sigma (Σ) hypernuclei have been sought, as have doubly-strange nuclei containing xi baryons (Ξ) or two Λ's.
Nomenclature
Hypernuclei are named in terms of their atomic number and baryon number, as in normal nuclei, plus the hyperon(s) which are listed in a left subscript of the symbol, with the caveat that atomic number is interpreted as the total charge of the hypernucleus, including charged hyperons such as the xi minus (Ξ−) as well as protons. For example, the hypernucleus contains 8 protons, 7 neutrons, and one Λ (which carries no charge).
History
The first was discovered by Marian Danysz and Jerzy Pniewski in 1952 using a nuclear emulsion plate exposed to cosmic rays, based on their energetic but delayed decay. This event was inferred to be due to a nuclear fragment containing a Λ baryon. Experiments until the 1970s would continue to study hypernuclei produced in emulsions using cosmic rays, and later using pion (π) and kaon (K) beams from particle accelerators.
Since the 1980s, more efficient production methods using pion and kaon beams have allowed further investigation at various accelerator facilities, including CERN, Brookhaven National Laboratory, KEK, DAφNE, and JPARC. In the 2010s, heavy ion experiments such as ALICE and STAR first allowed the production and measurement of light hypernuclei formed through hadronization from quark–gluon plasma.
Properties
Hypernuclear physics differs from that of normal nuclei because a hyperon is distinguishable from the four nucleon spin and isospin. That is, a single hyperon is not restricted by the Pauli exclusion principle, and can sink to the lowest energy level. As such, hypernuclei are often smaller and more tightly bound than normal nuclei; for example, the lithium hypernucleus is 19% smaller than the normal nucleus 6Li. However, the hyperons can decay via the weak force; the mean lifetime of a free Λ is , and that of a Λ hypernucleus is usually slightly shorter.
A generalized mass formula developed for both the non-strange normal nuclei and strange hypernuclei can estimate masses of hypernuclei containing Λ, ΛΛ, Σ, and Ξ hyperon(s). The neutron and proton driplines for hypernuclei are predicted and existence of some exotic hypernuclei beyond the normal neutron and proton driplines are suggested. This generalized mass formula was named the "Samanta formula" by Botvina and Pochodzalla and used to predict relative yields of hypernuclei in heavy-ion collisions.
Types
Λ hypernuclei
The simplest, and most well understood, type of hypernucleus includes only the lightest hyperon, the Λ.
While two nucleons can interact through the nuclear force mediated by a virtual pion, the Λ becomes a Σ baryon upon emitting a pion, so the Λ–nucleon interaction is mediated solely by more massive mesons such as the η and ω mesons, or through the simultaneous exchange of two or more mesons. This means that the Λ–nucleon interaction is weaker and has a shorter range than the standard nuclear force, and the potential well of a Λ in the nucleus is shallower than that of a nucleon; in hypernuclei, the depth of the Λ potential is approximately 30 MeV. However, one-pion exchange in the Λ–nucleon interaction does cause quantum-mechanical mixing of the Λ and Σ baryons in hypernuclei (which does not happen in free space), especially in neutron-rich hypernuclei. Additionally, the three-body force between a Λ and two nucleons is expected to be more important than the three-body interaction in nuclei, since the Λ can exchange two pions with a virtual Σ intermediate, while the equivalent process in nucleons requires a relatively heavy delta baryon (Δ) intermediate.
Like all hyperons, Λ hypernuclei can decay through the weak interaction, which changes it to a lighter baryon and emits a meson or a lepton–antilepton pair. In free space, the Λ usually decays via the weak force to a proton and a π– meson, or a neutron and a π0, with a total half-life of . A nucleon in the hypernucleus can cause the Λ to decay via the weak force without emitting a pion; this process becomes dominant in heavy hypernuclei, due to suppression of the pion-emitting decay mode. The half-life of the Λ in a hypernucleus is considerably shorter, plateauing to about near , but some empirical measurements substantially disagree with each other or with theoretical predictions.
Hypertriton
The simplest hypernucleus is the hypertriton (), which consists of one proton, one neutron, and one Λ hyperon. The Λ in this system is very loosely bound, having a separation energy of 130 keV and a large radius of 10.6 fm, compared to about for the deuteron.
This loose binding would imply a lifetime similar to a free Λ. However, the measured hypertriton lifetime averaged across all experiments (about ) is substantially shorter than predicted by theory, as the non-mesonic decay mode is expected to be relatively minor; some experimental results are substantially shorter or longer than this average.
Σ hypernuclei
The existence of hypernuclei containing a Σ baryon is less clear. Several experiments in the early 1980s reported bound hypernuclear states above the Λ separation energy and presumed to contain one of the slightly heavier Σ baryons, but experiments later in the decade ruled out the existence of such states. Results from exotic atoms containing a Σ− bound to a nucleus by the electromagnetic force have found a net repulsive Σ–nucleon interaction in medium-sized and large hypernuclei, which means that no Σ hypernuclei exist in such mass range. However, an experiment in 1998 definitively observed the light Σ hypernucleus .
ΛΛ and Ξ hypernuclei
Hypernuclei containing two Λ baryons have been made. However, such hypernuclei are much harder to produce due to containing two strange quarks, and , only seven candidate ΛΛ hypernuclei have been observed. Like the Λ–nucleon interaction, empirical and theoretical models predict that the Λ–Λ interaction is mildly attractive.
Hypernuclei containing a Ξ baryon are known. Empirical studies and theoretical models indicate that the Ξ––proton interaction is attractive, but weaker than the Λ–nucleon interaction. Like the Σ– and other negatively charged particles, the Ξ– can also form an exotic atom. When a Ξ– is bound in an exotic atom or a hypernucleus, it quickly decays to a ΛΛ hypernucleus or to two Λ hypernuclei by exchanging a strange quark with a proton, which releases about 29 MeV of energy in free space:
Ξ− + p → Λ + Λ
Ω hypernuclei
Hypernuclei containing the omega baryon (Ω) were predicted using lattice QCD in 2018; in particular, the proton–Ω and Ω–Ω dibaryons (bound systems containing two baryons) are expected to be stable. , no such hypernuclei have been observed under any conditions, but the lightest such species could be produced in heavy-ion collisions, and measurements by the STAR experiment are consistent with the existence of the proton–Ω dibaryon.
Hypernuclei with higher strangeness
Since the Λ is electrically neutral and its nuclear force interactions are attractive, there are predicted to be arbitrarily large hypernuclei with high strangeness and small net charge, including species with no nucleons. Binding energy per baryon in multi-strange hypernuclei can reach up to 21 MeV/A under certain conditions, compared to 8.80 MeV/A for the ordinary nucleus 62Ni. Additionally, formation of Ξ baryons should quickly become energetically favorable, unlike when there are no Λ's, because the exchange of strangeness with a nucleon would be impossible due to the Pauli exclusion principle.
Production
Several modes of production have been devised to make hypernuclei through bombardment of normal nuclei.
Strangeness exchange and production
One method of producing a K− meson exchanges a strange quark with a nucleon and changes it to a Λ:
p + K− → Λ + π0
n + K− → Λ + π−
The cross section for the formation of a hypernucleus is maximized when the momentum of the kaon beam is approximately 500 MeV/c. Several variants of this setup exist, including ones where the incident kaons are either brought to rest before colliding with a nucleus.
In rare cases, the incoming K− can instead produce a Ξ hypernucleus via the reaction:
p + K− → Ξ− + K+
The equivalent strangeness production reaction involves a π+ meson reacts with a neutron to change it to a Λ:
n + π+ → Λ + K+
This reaction has a maximum cross section at a beam momentum of 1.05 GeV/c, and is the most efficient production route for Λ hypernuclei, but requires larger targets than strangeness exchange methods.
Elastic scattering
Electron scattering off of a proton can change it to a Λ and produce a K+:
p + e− → Λ + e− + K+
where the prime symbol denotes a scattered electron. The energy of an electron beam can be more easily tuned than pion or kaon beams, making it easier to measure and calibrate hypernuclear energy levels. Initially theoretically predicted in the 1980s, this method was first used experimentally in the early 2000s.
Hyperon capture
The capture of a Ξ− baryon by a nucleus can make a Ξ− exotic atom or hypernucleus. Upon capture, it changes to a ΛΛ hypernucleus or two Λ hypernuclei. The disadvantage is that the Ξ− baryon is harder to make into a beam than singly strange hadrons. However, an experiment at J-PARC begun in 2020 will compile data on Ξ and ΛΛ hypernuclei using a similar, non-beam setup where scattered Ξ− baryons rain onto an emulsion target.
Heavy-ion collisions
Similar species
Kaonic nuclei
The K– meson can orbit a nucleus in an exotic atom, such as in kaonic hydrogen. Although the K–-proton strong interaction in kaonic hydrogen is repulsive, the K––nucleus interaction is attractive for larger systems, so this meson can enter a strongly bound state closely related to a hypernucleus; in particular, the K––proton–proton system is experimentally known and more tightly bound than a normal nucleus.
Charmed hypernuclei
Nuclei containing a charm quark have been predicted theoretically since 1977, and are described as charmed hypernuclei despite the possible absence of strange quarks. In particular, the lightest charmed baryons, the Λc and Σc baryons, are predicted to exist in bound states in charmed hypernuclei, and could be created in processes analogous to those used to make hypernuclei. The depth of the Λc potential in nuclear matter is predicted to be 58 MeV, but unlike Λ hypernuclei, larger hypernuclei containing the positively charged Λc would be less stable than the corresponding Λ hypernuclei due to Coulomb repulsion. The mass difference between the Λc and the is too large for appreciable mixing of these baryons to occur in hypernuclei. Weak decays of charmed hypernuclei have strong relativistic corrections compared to those in ordinary hypernuclei, as the energy released in the decay process is comparable to the mass of the Λ baryon.
Antihypernuclei
In August 2024 the STAR Collaboration reported the observation of the heaviest antimatter nucleus known, antihyperhydrogen-4 consisting of one antiproton, two antineutrons and an antihyperon.
The anti-lambda hyperon and the antihypertriton have also been previously observed.
| Physical sciences | Nuclear physics | Physics |
344536 | https://en.wikipedia.org/wiki/Hepatitis%20A | Hepatitis A | Hepatitis A is an infectious disease of the liver caused by Hepatovirus A (HAV); it is a type of viral hepatitis. Many cases have few or no symptoms, especially in the young. The time between infection and symptoms, in those who develop them, is two–six weeks. When symptoms occur, they typically last eight weeks and may include nausea, vomiting, diarrhea, jaundice, fever, and abdominal pain. Around 10–15% of people experience a recurrence of symptoms during the 6 months after the initial infection. Acute liver failure may rarely occur, with this being more common in the elderly.
It is usually spread by eating food or drinking water contaminated with infected feces. Undercooked or raw shellfish are relatively common sources. It may also be spread through close contact with an infectious person. While children often do not have symptoms when infected, they are still able to infect others. After a single infection, a person is immune for the rest of their life. Diagnosis requires blood testing, as the symptoms are similar to those of a number of other diseases. It is one of five known hepatitis viruses: A, B, C, D, and E.
The hepatitis A vaccine is effective for prevention.
Some countries recommend it routinely for children and those at higher risk who have not previously been vaccinated. It appears to be effective for life. Other preventive measures include hand washing and properly cooking food. No specific treatment is available, with rest and medications for nausea or diarrhea recommended on an as-needed basis. Infections usually resolve completely and without ongoing liver disease. Treatment of acute liver failure, if it occurs, is with liver transplantation.
Globally, around 1.4 million symptomatic cases occur each year and about 114 million infections (symptomatic and asymptomatic). It is more common in regions of the world with poor sanitation and not enough safe water. In the developing world, about 90% of children have been infected by age 10, thus are immune by adulthood. It often occurs in outbreaks in moderately developed countries where children are not exposed when young and vaccination is not widespread. Acute hepatitis A resulted in 11,200 deaths in 2015. World Hepatitis Day occurs each year on July 28 to bring awareness to viral hepatitis.
Signs and symptoms
Early symptoms of hepatitis A infection can be mistaken for influenza, but some people, especially children, exhibit no symptoms at all. Symptoms typically appear two–six weeks (the incubation period) after the initial infection. About 90% of children do not have symptoms. The time between infection and symptoms, in those who develop them, is two–six weeks, with an average of 28 days.
The risk for symptomatic infection is directly related to age, with more than 80% of adults having symptoms compatible with acute viral hepatitis and the majority of children having either asymptomatic or unrecognized infections.
Symptoms usually last less than 2 months, although some people can be ill for as long as 6 months:
Fatigue
Fever
Nausea
Appetite loss
Jaundice, a yellowing of the skin or the whites of the eyes owing to hyperbilirubinemia
Bile is removed from the bloodstream and excreted in the urine, giving it a dark amber color
Diarrhea
Light or clay-colored faeces (acholic faeces)
Abdominal discomfort
Extrahepatic manifestations
Joint pains, red cell aplasia, pancreatitis and generalized lymphadenopathy are the possible extrahepatic manifestations. Kidney failure and pericarditis are very uncommon. If they occur, they show an acute onset and disappear upon resolution of the disease.
Virology
Taxonomy
Hepatovirus A is a species of virus in the order Picornavirales, family Picornaviridae, genus Hepatovirus. Humans and other vertebrates serve as natural hosts of this genus.
Nine members of Hepatovirus are recognized. These species infect bats, rodents, hedgehogs, and shrews. Phylogenetic analysis suggests a rodent origin for human Hepatitis A.
A member virus of hepatovirus B (Phopivirus) has been isolated from a seal. This virus shared a common ancestor with Hepatovirus A about 1800 years ago.
Another hepatovirus – Marmota himalayana hepatovirus – has been isolated from the woodchuck Marmota himalayana. This virus appears to have had a common ancestor with the primate-infecting species around 1000 years ago.
Genotypes
One serotype and six different genotypes (three human and three simian) have been described. The human genotypes are numbered I–III. Six subtypes have been described (IA, IB, IIA, IIB, IIIA, IIIB). The simian genotypes have been numbered IV–VI. A single isolate of genotype VII isolated from a human has also been described but has been reclassified as subgenotype IIB. Genotype III has been isolated from both humans and owl monkeys. Most human isolates are of genotype I. Of genotype I isolates, subtype IA accounts for the majority.
The mutation rate in the genome has been estimated to be nucleotide substitutions per site per year. The human strains appear to have diverged from the simian about 3600 years ago. The mean age of genotypes III and IIIA strains has been estimated to be 592 and 202 years, respectively.
Structure
Hepatovirus A is a picornavirus; it is not enveloped and contains a positive-sense, single-strand of RNA packaged in a protein shell. Only one serotype of the virus has been found, but multiple genotypes exist. Codon use within the genome is biased and unusually distinct from its host. It also has a poor internal ribosome entry site. In the region that codes for the HAV capsid, highly conserved clusters of rare codons restrict antigenic variability.
Replication cycle
Vertebrates such as humans serve as the natural hosts. Transmission routes are fecal-oral and blood.
Following ingestion, HAV enters the bloodstream through the epithelium of the oropharynx or intestine. The blood carries the virus to its target, the liver, where it multiplies within hepatocytes and Kupffer cells (liver macrophages). Viral replication is cytoplasmic. Entry into the host cell is achieved by attachment of the virus to host receptors, which mediates endocytosis. Replication follows the positive-stranded RNA virus replication model. Translation takes place by viral initiation. The virus exits the host cell by lysis and viroporins. Virions are secreted into the bile and released in stool. HAV is excreted in large numbers about 11 days prior to the appearance of symptoms or anti-HAV IgM antibodies in the blood. The incubation period is 15–50 days and risk of death in those infected is less than 0.5%.
Within the liver hepatocytes, the RNA genome is released from the protein coat and is translated by the cell's own ribosomes. Unlike other picornaviruses, this virus requires an intact eukaryotic initiation factor 4G (eIF4G) for the initiation of translation. The requirement for this factor results in an inability to shut down host protein synthesis, unlike other picornaviruses. The virus must then inefficiently compete for the cellular translational machinery, which may explain its poor growth in cell culture. Aragonès et al. (2010) theorize that the virus has evolved a naturally highly deoptimized codon usage with respect to that of its cellular host in order to negatively influence viral protein translation kinetics and allow time for capsid proteins to fold optimally.
No apparent virus-mediated cytotoxicity occurs, presumably because of the virus' own requirement for an intact eIF4G, and liver pathology is likely immune-mediated.
Transmission
The virus primarily spreads by the fecal–oral route, and infections often occur in conditions of poor sanitation and overcrowding. Hepatitis A can be transmitted by the parenteral route, but very rarely by blood and blood products. Food-borne outbreaks are common, and ingestion of shellfish cultivated in polluted water is associated with a high risk of infection. HAV can also be spread through sexual contact, specifically oro–anal and digital–rectal sexual acts. Humans are the only natural reservoir and disease vector of the HAV virus; no known insect or other animal vectors can transmit the virus. A chronic HAV state has not been reported.
About 40% of all acute viral hepatitis is caused by HAV. Infected individuals are infectious prior to onset of symptoms, roughly 10 days following infection. The virus is resistant to detergent, acid (pH 1), solvents (e.g., ether, chloroform), drying, and temperatures up to 60 °C. It can survive for months in fresh and salt water. Common-source (e.g., water, food) outbreaks are typical. Infection is common in children in developing countries, reaching 100% incidence, but following infection, lifelong immunity results. HAV can be inactivated by chlorine treatment (drinking water), formalin (0.35%, 37 °C, 72 hours), peracetic acid (2%, 4 hours), beta-propiolactone (0.25%, 1 hour), and UV radiation (2 μW/cm2/min).
In developing countries, and in regions with poor hygiene standards, the rates of infection with this virus are high and the illness is usually contracted in early childhood. As incomes rise and access to clean water increases, the incidence of HAV decreases. In developed countries, though, the infection is contracted primarily by susceptible young adults, most of whom are infected with the virus during trips to countries with a high incidence of the disease or through contact with infectious persons.
Diagnosis
Although HAV is excreted in the feces towards the end of the incubation period, specific diagnosis is made by the detection of HAV-specific IgM antibodies in the blood. IgM antibody is only present in the blood following an acute hepatitis A infection. It is detectable from 1–2 weeks after the initial infection and persists for up to 14 weeks. The presence of IgG antibodies in the blood means the acute stage of the illness has passed and the person is immune to further infection. IgG antibodies to HAV are also found in the blood following vaccination, and tests for immunity to the virus are based on the detection of these antibodies.
During the acute stage of the infection, the liver enzyme alanine transferase (ALT) is present in the blood at levels much higher than is normal. The enzyme comes from the liver cells damaged by the virus.
Hepatovirus A is present in the blood (viremia) and feces of infected people up to 2 weeks before clinical illness develops.
Prevention
Hepatitis A can be prevented by vaccination, good hygiene, and sanitation.
Vaccination
The two types of vaccines contain either inactivated Hepatovirus A or a live but attenuated virus. Both provide active immunity against a future infection. The vaccine protects against HAV in more than 95% of cases for longer than 25 years. In the United States, the vaccine developed by Maurice Hilleman and his team was licensed in 1995, and the vaccine was first used in 1996 for children in high-risk areas, and in 1999 it was spread to areas with elevating levels of infection.
The vaccine is given by injection. An initial dose provides protection lasting one year starting 2–4 weeks after vaccination; the second booster dose, given six to 12 months later, provides protection for over 20 years.
The vaccine was introduced in 1992 and was initially recommended for persons at high risk. Since then, Bahrain and Israel have embarked on elimination programmes. In countries where widespread vaccination has been practised, the incidence of hepatitis A has decreased dramatically.
In the United States, vaccination of children is recommended at 1 and 2 years of age; hepatitis A vaccination is not recommended in those younger than 12 months of age. It is also recommended in those who have not been previously immunized and who have been exposed or are likely to be exposed due to travel. The CDC recommends vaccination against infection for men who have sex with men.
Treatment
No specific treatment for hepatitis A is known. Recovery from symptoms following infection may take several weeks or months. Therapy is aimed at maintaining comfort and adequate nutritional balance, including replacement of fluids lost from vomiting and diarrhea.
Prognosis
In the United States in 1991, the mortality rate for hepatitis A was estimated to be 0.015% for the general population, but ranged up to 1.8–2.1% for those aged 50 and over who were hospitalized with icteric hepatitis. The risk of death from acute liver failure following HAV infection increases with age and when the person has underlying chronic liver disease.
Young children who are infected with hepatitis A typically have a milder form of the disease, usually lasting 1–3 weeks, whereas adults tend to experience a much more severe form of the disease.
Epidemiology
Globally, symptomatic HAV infections are believed to occur in around 1.4 million people a year. About 114 million infections (asymptomatic and symptomatic) occurred all together in 2015. Acute hepatitis A resulted in 11,200 deaths in 2015. Developed countries have low circulating levels of hepatovirus A, while developing countries have higher levels of circulation. Most adolescents and adults in developing countries have already had the disease, thus are immune. Adults in midlevel countries may be at risk of disease with the potential of being exposed.
Countries
Over 30,000 cases of hepatitis A were reported to the CDC in the US in 1997, but the number has since dropped to less than 2,000 cases reported per year.
The most widespread hepatitis A outbreak in the United States occurred in 2018, in the state of Kentucky. The outbreak is believed to have started in November 2017. By July 2018 48% of the state's counties had reported at least one case of hepatitis A, and the total number of suspected cases was 969 with six deaths (482 cases in Louisville, Kentucky). By July 2019 the outbreak had reached 5,000 cases and 60 deaths, but had slowed to just a few new cases per month.
Another widespread outbreak in the United States, the 2003 US hepatitis outbreak, affected at least 640 people (killing four) in northeastern Ohio and southwestern Pennsylvania in late 2003. The outbreak was blamed on tainted green onions at a restaurant in Monaca, Pennsylvania. In 1988, more than 300,000 people in Shanghai, China, were infected with HAV after eating clams (Anadara subcrenata) from a contaminated river.
In June 2013, frozen berries sold by US retailer Costco and purchased by around 240,000 people were the subject of a recall, after at least 158 people were infected with HAV, 69 of whom were hospitalized. In April 2016, frozen berries sold by Costco were once again the subject of a recall, after at least 13 people in Canada were infected with HAV, three of whom were hospitalized. In Australia in February 2015, a recall of frozen berries was issued after at least 19 people contracted the illness following their consumption of the product. In 2017, California (particularly around San Diego), Michigan, and Utah reported outbreaks of hepatitis A that have led to over 800 hospitalizations and 40 deaths.
| Biology and health sciences | Viral diseases | Health |
344618 | https://en.wikipedia.org/wiki/China%20National%20Space%20Administration | China National Space Administration | The China National Space Administration (CNSA) is a government agency of the People's Republic of China headquartered in Haidian, Beijing, responsible for civil space administration and international space cooperation. These responsibilities include organizing or leading foreign exchanges and cooperation in the aerospace field. The CNSA is an administrative agency under the Ministry of Industry and Information Technology.
Founded in 1993, CNSA has pioneered a number of achievements in space for China despite its relatively short history, including becoming the first space agency to land on the far side of the Moon with Chang'e 4, bringing material back from the Moon with Chang'e 5 and 6, and being the second agency who successfully landed a rover on Mars with Tianwen-1.
As the governing body of civil space activities, China National Space Administration does not execute any space program. The China Aerospace Science and Technology Corporation executes China's state space programs instead. The China Manned Space Program is operated by China Manned Space Agency, instead of the CNSA.
History
CNSA is an agency created in 1993 when the Ministry of Aerospace Industry was split into CNSA and the China Aerospace Science and Technology Corporation (CASC). The former was to be responsible for policy, while the latter was to be responsible for execution. This arrangement proved somewhat unsatisfactory, as these two agencies were, in effect, one large agency, sharing both personnel and management.
As part of a massive restructuring in 1998, CASC was split into a number of smaller state-owned companies. The intention appeared to have been to create a system similar to that characteristic of Western defense procurement in which entities which are government agencies, setting operational policy, would then contract out their operational requirements to entities which were government-owned, but not government-managed.
Since the passage of the Wolf Amendment in 2011, NASA has been forced by Congress to implement a long-standing exclusion policy with CNSA ever since, though this has been periodically overcome.
In 2021, China began building the Tiangong space station, which consists of three modules designated for crew, cargo, and research. The construction was completed in late 2022, and there are plans to add an additional three modules.
In 2024, China announced that it will undertake 100 space missions, a significant increase from the 70 missions conducted in 2023 this is mostly satellites, testing, crew replacement, cargo, and more.
Function
CNSA was established as a government institution to develop and fulfill China's due international obligations, with the approval by the 8th National People's Congress of China (NPC). The 9th NPC assigned CNSA as an internal structure of the Commission of Science, Technology and Industry for National Defense (COSTIND).
CNSA assumes the following main responsibilities: signing governmental agreements in the space area on behalf of organizations, inter-governmental scientific and technical exchanges; and also being in charge of the enforcement of national space policies and managing the national space science, technology and industry.
China has signed governmental space cooperation agreements with Argentina, Brazil, Chile, France, Germany, India, Italy, Pakistan, Russia, Ukraine, the United Kingdom, the United States, and some other countries. Significant achievements have been scored in the bilateral and multilateral and technology exchanges and cooperation.
Administrators
The most recent administrator is Zhang Kejian. Wu Yanhua is vice-administrator and Tian Yulong is secretary general.
April 1993: Liu Jiyuan
April 1998: Luan Enjie
2004: Sun Laiyan
July 2010: Chen Qiufa
March 2013: Ma Xingrui
December 2013: Xu Dazhe
May 2017: Tang Dengjie
May 2018: Zhang Kejian
Departments
There are four departments under the CNSA:
Department of General Planning
Department of System Engineering
Department of Science, Technology and Quality Control
Department of Foreign Affairs
CNSA's logo is a similar design to that of China Aerospace Science and Technology Corporation. The arrow in the middle is similar to the Chinese character 人 which means 'human' or 'people', to state that humans are the center of all space exploration. The three concentric ellipses stand for three types of escape velocity (minimum speed needed to reach sustainable orbits, to escape the Earth system, and to escape the Solar System) which are milestones of space exploration. The second ring is drawn with a bold line, to state that China has passed the first stage of exploration (Earth system) and is undergoing the second stage exploration (within the Solar System). The 人 character stands above the three rings to emphasize humanity's capability to escape and explore. Olive branches were added to state that China's space exploration is peaceful in nature.
Launch facilities
Jiuquan Satellite Launch Center
Taiyuan Satellite Launch Center
Xichang Satellite Launch Center
Wenchang Space Launch Site
| Technology | Programs and launch sites | null |
344673 | https://en.wikipedia.org/wiki/Apocynaceae | Apocynaceae | Apocynaceae (, from Apocynum, Greek for "dog-away") is a family of flowering plants that includes trees, shrubs, herbs, stem succulents, and vines, commonly known as the dogbane family, because some taxa were used as dog poison. Members of the family are native to the European, Asian, African, Australian, and American tropics or subtropics, with some temperate members. The former family Asclepiadaceae (now known as Asclepiadoideae) is considered a subfamily of Apocynaceae and contains 348 genera. A list of Apocynaceae genera may be found here.
Many species are tall trees found in tropical forests, but some grow in tropical dry (xeric) environments. Also perennial herbs from temperate zones occur. Many of these plants have milky latex, and many species are poisonous if ingested, the family being rich in genera containing alkaloids and cardiac glycosides, those containing the latter often finding use as arrow poisons. Some genera of Apocynaceae, such as Adenium, bleed clear sap without latex when damaged, and others, such as Pachypodium, have milky latex apart from their sap.
Description
Growth pattern
The dogbane/milkweed family includes annual plants, perennial herbs, stem succulents, woody shrubs, trees, or vines. Most exude a milky latex when cut.
Leaves and stems
Leaves are simple. They may appear one at a time (singly) with each occurrence on alternating sides of the stem, but usually occur in pairs (and rarely in whorls). When paired, they occur on opposite sides of the stem (opposite), with each pair occurring at an angle rotated 90° to the pair below it (decussate).
There is no stipule (a small leaf-like structure at the base of the leaf stem), or stipules are small and sometimes finger-like.
Inflorescence and fruit
Flowers have radial symmetry (actinomorphic), and are borne in heads that are cymes or racemes, or are solitary in axils. They are perfect (bisexual), with a synsepalous, five-lobed calyx united into a tube at the base. Inflorescences are terminal or axillary. Five petals are united into a tube with four or five epipetalous stamens. The style head is swollen. The pollen is transported in foam. The ovary is usually superior, bicarpellary, and apocarpous, with a common fused style and stigma. (Fig. 5. and Fig.6. in the illustration of Rhigospira quadrangularis show a typical tripartite style which divides into three zones (specialised for pollen deposition, viscin secretion, and the reception of pollen).
The fruit is a drupe, a berry, a capsule, or a (frequently paired) follicle. The seeds are often winged or have appendages of long silky hairs.
Taxonomy
As of 2012, the family was described as comprising some 5,100 species, in five subfamilies:
Apocynoideae Burnett, 1835
Asclepiadoideae Burnett, 1835 (incorporating the Asclepiadaceae)
Periplocoideae Endl., 1838
Rauvolfioideae Kostel., 1834
Secamonoideae Endl., 1838
The former family Asclepiadaceae is included in Apocynaceae according to the Angiosperm Phylogeny Group III (APG III) modern, largely molecular-based system of flowering plant taxonomy.
An updated classification, including 366 genera, 25 tribes, and 49 subtribes, was published in 2014.
376 genera are currently accepted.
Distribution and habitat
Species in this family are distributed mainly in tropical regions:
In the tropical forests and swamps of Indomalaya: small to very tall evergreen trees up to tall, often with buttress roots, such as Alstonia and Dyera.
In Australia: occurs in all habitats; about 46 genera and about 200 species, including about 20 naturalised; herbs, vines, shrubs and trees.
In deciduous forests of Africa, India, and Indo-China: smaller trees such as Carissa, Wrightia, and Holarrhena
In tropical America, India, Myanmar, and Malaya: evergreen trees and shrubs, such as Rauvolfia, Tabernaemontana, and Acokanthera.
In Central America: Plumeria, or the frangipani, with its waxy white or pink flowers and a sweet scent.
In South America, Africa, and Madagascar: many lianas, such as Landolphia
In the Mediterranean region: Nerium, with the well-known oleander or be-still tree (Nerium oleander), and Apple of Sodom (Calotropis procera), with other (Calotropis) species extending into South Asia.
The only genera found in temperate Europe away from the Mediterranean are Vinca (Rauvolfioideae) and Vincetoxicum (Asclepiadoideae). Also Asclepias syriaca is an invasive weed (e. g., in many areas of Ukraine).
In North America: Apocynum, dogbane or Indian hemp, including Apocynum cannabinum, a traditional source of fiber. Also the bluestars, Amsonia, herbaceous perennials of upright habit, grown as ornamental plants for their attractive flowers.
In continental southern Africa (Angola, Botswana, Eswatini, Mozambique, South Africa, and Zimbabwe) and Madagascar, except for the humid evergreen forest of the eastern side of Madagascar, and never above for the entire island: Pachypodium and Fockea.
Ecology
Several genera are preferred larval host plants for the Queen Butterfly (Danaus gilippus).
Toxicity
Many species of plants from the family Apocynaceae have some toxicity, with some being extremely poisonous if parts are ingested, or if they are not handled properly. Genera containing cardiac glycosides—Cerbera, Nerium, Asclepias, Cascabela, Strophanthus, Acokanthera, Apocynum, Thevetia, etc.—have therapeutic ranges, but are often associated with accidental poisonings, in many cases lethal (see below). Alkaloid-producing species like Rauvolfia serpentina, Catharanthus roseus, and Tabernanthe iboga are likewise the source of compounds with therapeutic ranges, but which have significant associated toxicities if not taken in appropriate doses and in controlled fashion. (See below)
Uses
Several members of the family Apocynaceae have had economic uses in the past. Several are sources of important natural products—pharmacologic tool compounds and drug research candidates, and in some cases actual prescription drugs. Cardiac glycosides, which affect heart function, are a ready example. Genera studied and known to contain such glycosides include Acokanthera, Apocynum, Cerbera, Nerium, Thevetia and Strophanthus. Rauvolfia serpentina (Indian snakeroot) contains the alkaloid reserpine, which has been used as an antihypertensive and an antipsychotic drug but its adverse effects limit its clinical use. Catharanthus roseus yields alkaloids used in the treatment of cancer. Tabernanthe iboga, Voacanga africana, and Tabernaemontana undulata contain the alkaloid ibogaine, which is a psychedelic drug which may help with drug addiction, but which has significant adverse effects, with ibogaine being both cardiotoxic and neurotoxic. Ajmalicine, an alkaloid found in Rauvolfia spp., Catharanthus roseus, and Mitragyna speciosa, is an antihypertensive drug used in the treatment of high blood pressure.
Many genera are grown as ornamental plants, including Amsonia (bluestar), Nerium (oleander), Vinca (periwinkle), Carissa (Natal plum), Allamanda (golden trumpet), Plumeria (frangipani), Thevetia, Mandevilla (Savannah flower), and Adenium (desert-rose).
In addition, the genera Landolphia, Carpodinus, and Mascarenhasia have been used as commercial sources of inferior rubber. (See Congo rubber)
There are limited dietary uses of plants from this family. The flower of Echites panduratus (common name: ) is edible. Carissa (Natal plum) produces an edible fruit, but all other parts of the plant are poisonous. The genus Apocynum was reportedly used as a source of fiber by Native Americans. The aromatic fruit juice from Saba comorensis (syn. Landolphia comorensis, the Bungo or Mbungo fruit) is used as a drink.
Finally, ethnopharmacologic and ethnotoxicologic uses are also known. The roots of Tabernanthe iboga and certain Voacanga species have traditionally been used ceremonially as hallucinogens in Africa. The ibogaine-type alkaloids responsible for the psychoactivity of these plants have been studied with regard to the treatment of drug addiction. The juice of Acokanthera species such as A. venenata and the milky juice of the Namibian Pachypodium have been used as poison for arrow tips.
Many species are ornamental in gardens or as houseplants.
Gallery
Flowers
Fruits
Pachycaul species
| Biology and health sciences | Others | null |
344775 | https://en.wikipedia.org/wiki/Earthenware | Earthenware | Earthenware is glazed or unglazed nonvitreous pottery that has normally been fired below . Basic earthenware, often called terracotta, absorbs liquids such as water. However, earthenware can be made impervious to liquids by coating it with a ceramic glaze, and such a process is used for the great majority of modern domestic earthenware. The main other important types of pottery are porcelain, bone china, and stoneware, all fired at high enough temperatures to vitrify. End applications include tableware and decorative ware such as figurines.
Earthenware comprises "most building bricks, nearly all European pottery up to the seventeenth century, most of the wares of Egypt, Persia and the near East; Greek, Roman and Mediterranean, and some of the Chinese; and the fine earthenware which forms the greater part of our tableware today" ("today" being 1962). Pit fired earthenware dates back to as early as 29,000–25,000 BC, and for millennia, only earthenware pottery was made, with stoneware gradually developing some 5,000 years ago, but then apparently disappearing for a few thousand years. Outside East Asia, porcelain was manufactured at any scale only from the 18th century AD, and then initially as an expensive luxury.
After it is fired, earthenware is opaque and non-vitreous, soft and capable of being scratched with a knife. The Combined Nomenclature of the European Union describes it as being made of selected clays sometimes mixed with feldspars and varying amounts of other minerals, and white or light-coloured (i.e., slightly greyish, cream, or ivory).
Characteristics
Generally, unfired earthenware bodies exhibit higher plasticity than most whiteware bodies and hence are easier to shape by RAM press, roller-head or potter's wheel than bone china or porcelain.
Due to its porosity, fired earthenware, with a water absorption of 5-8%, must be glazed to be watertight. Earthenware has lower mechanical strength than bone china, porcelain or stoneware, and consequently articles are commonly made in thicker cross-section, although they are still more easily chipped.
Darker-coloured terracotta earthenware, typically orange or red due to a comparatively high content of iron oxides, are widely used for flower pots, tiles and some decorative and oven ware.
Production
Materials
The compositions of earthenware bodies vary considerably, and include both prepared and 'as dug'; the former being by far the dominant type for studio and industry. A general body formulation for contemporary earthenware is 25% kaolin, 25% ball clay, 35% quartz and 15% feldspar.
Shaping
Firing
Earthenware can be produced at firing temperatures as low as and many clays will not fire successfully above about . Much historical pottery was fired somewhere around , giving a wide margin of error where there was no precise way of measuring temperature, and very variable conditions within the kiln.
Modern earthenware may be biscuit (or "bisque") fired to temperatures between and glost-fired (or "glaze-fired") to between . Some studio potters follow the reverse practice, with a low-temperature biscuit firing and a high-temperature glost firing. Oxidising atmospheres are the most common.
After firing, most earthenware bodies will be colored white, buff or red. For iron-rich bodies earthenware, firing at comparatively low temperature in an oxidising atmosphere results in a red colour, whilst higher temperatures with a reducing atmosphere results in darker colours, including black. Higher firing temperatures may cause earthenware to bloat.
Examples of earthenware
Despite the most highly valued types of pottery often switching to stoneware and porcelain as these were developed by a particular culture, there are many artistically important types of earthenware. All ancient Greek and ancient Roman pottery is earthenware, as is the Hispano-Moresque ware of the late Middle Ages, which developed into tin-glazed pottery or faience traditions in several parts of Europe, mostly notably the painted maiolica of the Italian Renaissance, and Dutch Delftware. With a white glaze, these were able to imitate porcelains both from East Asia and Europe.
Amongst the most complicated earthenware ever made are the life-size Yixian glazed pottery luohans of the Liao dynasty (907–1125), Saint-Porchaire ware of the mid-16th century, apparently made for the French court and the life-size majolica peacocks by Mintons in the 1860s.
In the 18th century, especially in English Staffordshire pottery, technical improvements enabled very fine wares such as Wedgwood's creamware, that competed with porcelain with considerable success, as his huge creamware Frog Service for Catherine the Great showed. The invention of transfer printing processes made highly decorated wares cheap enough for far wider sections of the population in Europe.
In China, sancai glazed wares were lead-glazed earthenware, and as elsewhere, terracotta remained important for sculpture. The Etruscans had made large sculptures such as statues in it, where the Romans used it mainly for figurines and Campana reliefs. Chinese painted or Tang dynasty tomb figures were earthenware as were the later Yixian glazed pottery luohans. After the ceramic figurine was revived in European porcelain, earthenware figures followed, such as the popular English Staffordshire figures.
| Technology | Materials | null |
344899 | https://en.wikipedia.org/wiki/Quercus%20suber | Quercus suber | Quercus suber, commonly called the cork oak, is a medium-sized, evergreen oak tree in the section Quercus sect. Cerris. It is the primary source of cork for wine bottle stoppers and other uses, such as cork flooring and as the cores of cricket balls. It is native to southwest Europe and northwest Africa. In the Mediterranean basin the tree is an ancient species with fossil remnants dating back to the Tertiary period. It can survive for as long as two centuries. Typically, once it reaches 25 years old, its thick bark can be harvested for cork every 9 to 12 years without causing harm to the tree.
It endures drought and makes little demand on the soil quality and is regarded as a defence against desertification. Cork oak woodlands are home to a multitude of animal and plant species. Since cork for sealing bottles is increasingly being displaced by other materials, these forests are at risk as part of the cultural landscape and as a result animal species such as the Iberian lynx and imperial eagles are threatened with extinction.
Description
General appearance and bark
The cork oak grows as an evergreen tree, reaching an average height of or in rare cases up to 25 m and a trunk diameter (DBH) of . It forms a dense and asymmetrical crown that starts at a height of and spreads widely in free-standing trees. The crown can be divided into several separate, rounded partial crowns.
The young twigs are densely hairy light gray or whitish. Older branches are strong and knotty. Older trees only form short shoots between in length.
The thick, longitudinally cracked cork layers of the gray-brown trunk bark are characteristic of the cork oak. The cambium of the smooth bark of young trees forms a cork layer very early on, which can be thick. The light and spongy cork fabric shows vertical cracks and is white on the outside and red to red-brown on the inside. After the cork has been harvested, the trunk appears reddish brown, but later it is significantly darker. The wood is ring-pored, has a brown heartwood and a light reddish sapwood. The cork oak develops a taproot that reaches a depth of and from which several meters long, horizontally running side roots extend. The trees can live over 400 years, and harvested specimens can be 150 to 200 years old.
Leaves
The leathery leaves are alternate and are long and wide. The shape varies between round, oval and lanceolate-oval. The leaf blade has five to seven sharp teeth on both edges and a pointed vegetation cone (apex). The midrib stands out clearly on the underside of the leaf, the first-order lateral nerves usually lead to the teeth of the leaf margin. The upper side of the leaf is light green, the underside of the leaf whitish and densely hairy. There is no hair on young trees. The leaf stalks are long and are also hairy. At the base of the petiole are two narrow, lanceolate, long and bright red stipules that fall off in the first year. The new leaves appear in April and May, when older leaves are also shed. They usually stay on the tree for two to three years, less often only one year, the latter especially in severe environmental conditions and on the northern border of the distribution area. Extremely cold winters can also lead to complete defoliation.
Inflorescence and flower
The cork oak is single sexed (monoecious), with both female and male flowers on one specimen. The female flowers form upright inflorescences in the leaf axils of young branches. These are formed from a hairy axis long with two to five separate flowers. The female flowers contain a small, hairy, four- to six-lobed flower envelope and three to four styles. The male catkins also arise on the leaf axils of young branches. They are bright red at the beginning and stand upright, older catkins are yellow and pendulous, long and have a whitish hairy axis. The single flowers are sessile and have a densely hairy flower cover that is colored red when opened. The four to six stamens are whitish with yellow, egg-shaped anthers. They are longer than the bracts.
Infructescence, fruit and seed
The fruit clusters are long and carry two to eight acorns. About half of the fruits are enclosed in the fruit cup (cupule); the fruit cups are in diameter. The upper scales of the cupula are gray and hairy, in the subspecies Quercus suber occidentalis the scales are close together or are fused. The size of the acorns varies between lengths of and diameters of . The fruit casing (pericarp) is bare, smooth and shiny brownish red. The hilum (the starting point of the seed) is convex and has a diameter of .
Taxonomy
Quercus suber is a species of the section Cerris to which, for example, the following species also belong:
Valonia oak (Quercus macrolepis)
Turkey oak (Quercus cerris)
Quercus × crenata
Macedonian oak (Quercus trojana)
Characteristic for the section are the hairless pericarp and the usually two-year ripening time of the fruits. The cork oak is an exception because the fruits can ripen in both the first and the second year.
In the species Quercus suber two subspecies are distinguished:
Quercus suber subsp. suber: Nominal taxon
Quercus suber subsp. occidentalis (Gay) Bonnier & Layens: It differs from the nominate form in the shape of the cupula scales, the longer development time of the fruits and the semi-evergreen foliage. The distribution area of the subspecies is the Portuguese Atlantic coast.
Together with the Turkey oak (Quercus cerris) and the holm oak (Quercus ilex), the cork oak forms hybrids.
The scientific name Quercus suber is derived from the Latin word quercus, which the Romans used to describe the pedunculate oak (Quercus robur). The specific epithet suber means in Latin cork oak and also cork.
Distribution and habitat
The cork oak occupies the area around the western Mediterranean basin. In Portugal, natural and cultivated stands cover an area of 750,000 hectares. There are natural populations of the nominate form at altitudes between above sea level, the subspecies occidentalis is found along the Atlantic coast. In Spain the occurrences remain mostly below , but rarely reach heights of . In Spain, cork oaks are common in the southern half of the country, as well as in the western and northeastern areas, but rare in central Spain. In Italy one finds natural occurrences along the Tyrrhenian Sea and in eastern Apulia on the Adriatic Sea. Also on the Adriatic is the cork oak on the Dalmatian coast. It is one of the most common forest trees in Sardinia. Natural and man-made occurrences exist in Africa on the Mediterranean coast of Tunisia, Algeria and Morocco and at altitudes up to , on the High Atlas up to . In its native range, cork oak forests cover approximately . Outside of its natural range, the cork oak is cultivated in the Crimea, the Caucasus, India and the Southwestern United States. The subspecies Quercus suber occidentalis also thrives in mild areas of England.
The species needs very little light and cannot survive in dense populations. It loves warmth, grows at annual mean temperatures of and can withstand maximum temperatures of up to . In the area of distribution, the temperature rarely falls below freezing point, but temperatures down to without damage and down to without major damage can be tolerated. The cork oak is not hardy in Central Europe. It endures drought and survives dry periods in summer by reducing its metabolism. An annual rainfall of is considered optimal, in cooler locations can be sufficient with enough humidity. Cork oaks have low soil demands and also grow in poor, dry or rocky locations. They rarely thrive on calcareous soils, but they are often found on crystalline slates, on gneiss, granite and sands. The acidity of the soil should be between pH 4.5 and 7.
The cork oak is considered a pyrophyte because it recovers quickly after forest fires as it is protected by the cork.
Ecology
The cork oak forest is one of the major plant communities of the Mediterranean woodlands and forests ecoregion. In natural populations, the cork oak grows together with the holm oaks (Quercus ilex, Quercus rotundifolia), the Portuguese oak (Quercus faginea), the Pyrenean oak (Quercus pyrenaica), Mirbeck's oak (Quercus canariensis), the maritime pine (Pinus pinaster), the stone pine (Pinus pinea), the strawberry tree (Arbutus unedo) and the olive tree (Olea europaea), in cooler locations also with the sweet chestnut (Castanea sativa). In addition to these tree species, the shrub-forming species include the Kermes oak (Quercus coccifera), the Lusitanian oak (Quercus lusitanica) the holly buckthorn (Rhamnus alaternus), species of the genus Phillyrea, the myrtle (Myrtus communis), the green heather (Erica scoparia), the common smilax (Smilax aspera) and the Montpellier cistus (Cistus monspeliensis) are often found together with the cork oak.
As a pyrophyte, this tree has a thick, insulating bark that makes it well adapted to forest fires. After a fire, many tree species regenerate from seeds (as, for example, the maritime pine) or re-sprout from the base of the tree (as, for example, the holm oak). The bark of the cork oak allows it to survive fires and then simply regrow branches to fill out the canopy. The quick regeneration of this oak makes it successful in the fire-adapted ecosystems of the Mediterranean biome.
Symbiosis
The cork oak enters into a mycorrhizal symbiosis with several types of fungus. The fine root system of the oak is in close contact with the mycelium of the fungus. The oak receives water and nutrient salts from the fungus in exchange for products of photosynthesis. Such a symbiosis exists among others with the following species:
Caesar's mushroom (Amanita caesarea)
Death cap (Amanita phalloides)
Panther cap (Amanita pantherina)
Gilbert's limbed lepidella (Amanita gilbertii)
Cep (Boletus edulis)
Russula rubra
Diseases and predators
Cork oak is relatively resistant to pathogens, but some diseases occur in the species. Leaf spot can be caused by the fungus Apiognomonia errabunda. Other fungi can cause leaf scorching, powdery mildew, rust, and cankers.
The most virulent cork oak pathogen may be Diplodia corticola, a sac fungus which causes sap-bleeding sunken canker wounds in the wood, withering of the leaves, and lesions on the acorns. The fungus Biscogniauxia mediterranea is becoming more common in cork oak forests. Its fruiting bodies appear as charcoal-black cankers. Both of these fungi are transmitted by the oak pinhole borer (Platypus cylindrus), a species of weevil.
The common water mould Phytophthora cinnamomi grows in the roots of the tree and has been known to devastate cork oak woodlands.
Several species of butterflies damage the cork oak, the most important being the spongy moth (Lymantria dispar). The species lays its eggs in the bark of the branches and trunks, and the caterpillars that hatch in spring are distributed in the crown and eat them bare. The bacterial species Bacillus thuringiensis is used as a biological plant protection agent against the spongy moth. Another pest is the green oak tortrix (Tortrix viridana), whose caterpillars eat flowers and young leaves and roll them up with thread to form typical coils. The lackey moth (Malacosoma neustria) also causes damage to the leaves, sticking its eggs to the bark of thin twigs in multiple rows, and also the brown-tail moth (Euproctis chrysorrhoea), whose caterpillars skeletonize the leaves and further damage the tree after overwintering in spring. A special cork pest is the jewel beetle Coraebus undatus, which lays its eggs in the cork tissue. Another harmful species of beetle is the great capricorn beetle (Cerambyx cerdo), whose larvae eat long corridors in the oak wood.
Unfavorable climatic conditions and fungal attack are made responsible for the weakening of trees and for crown damage. Such fungal parasites of weakness are Botryosphaeria stevensii, Biscogniauxia mediterranea, Endothiella gyrosa and representatives of the mold genus Fusarium. Drought and parasite infestation are also considered to be the cause of the weakness syndrome in parts of Spain and Portugal.
Uses
The cork oak is grown for the production of cork in several Mediterranean countries. The centers of cork production are in southern Portugal (accounting for 50% of the total production) and southern Spain, where low trees with large crowns and strong branches are grown in large areas, which provide the highest yield of cork. These mostly extensively managed habitats are called montados in Portugal and dehesas in Spain. They are considered to be extremely valuable from the point of view of biodiversity and cultural heritage.
The cork consists of dead, air-filled, thin-walled cells and contains cellulose and suberin. Cork is heat and sound insulating, the suberine gives it water-repellent properties. The cork layer is replicated by the cork-producing phellogen and can therefore be harvested repeatedly without damaging the tree too much. The first harvest usually takes place after about 25 years with a trunk diameter (DBH) of , though new techniques (such as better irrigation systems) could shorten it to only 8 to 10 years. The first cork layer is called "male cork" or "virgin cork", is still not very elastic and cracked and is only used for insulating mats. The second harvested cork (known as secundeira), has a more regular structure and is softer, but is still only used for insulation and in decorative objects. Only the following cork harvests deliver a higher quality cork, the "female cork", which can be used commercially in full. The best quality cork is obtained from the third and fourth harvest. Cork harvesting takes place every nine to twelve years when a layer thickness of is reached. Under favorable (warm) conditions, the harvest can take place every eight years, in North Africa every seven years. A cork oak can be harvested five to seventeen times in total. In order to minimize the damage to the trunk surface, harvesting can be carried out every three years, whereby only a third of the usable surface is removed. An important maintenance measure is pruning, which begins around the age of ten at a height of about . Some sources say an oak can provide around of cork over its lifespan, and one hectare around per year while others suggest a single tree can produce on average of cork per harvest, a comparatively higher value, as cork oaks can live more than 200 years in good conditions.
The cork is mainly used for the production of stoppers and corks, as well as for heat and sound insulation, cork
paper, badminton shuttlecocks, cricket balls, handles of fishing rods and hand tools, special devices for the space industry and for other technical applications (including composite materials, shoe soles, floor coverings). Bottle cork production accounts for around 70% of the added value in cork cultivation. Since natural corks are increasingly being replaced by plastic or sheet metal closures, there could be a significant decline in the cork oak population in southwestern Europe, which endangers the biodiversity in these areas.
The bark, which contains around twelve percent extractable tannin, is also used. In addition, the acorns are used as feed in extensive pig fattening (acorn fattening), such as for Iberian ham production; although the Holm Oak (Quercus ilex), is preferred for this due to its sweeter fruits. One cork oak tree can provide of acorns per year.
Cork oaks cannot legally be cut down in Portugal, except for forest management felling of old, unproductive trees, and, even in those cases, farmers need special permission from the Ministry of Agriculture.
Cork harvesting is done entirely without machinery, being dependent solely on human labour. Usually five people are required to harvest the tree's bark, using a small axe. The process mandates specialized training due to the skill required to harvest bark without inflicting too much damage to the tree.
The European cork industry produces 300,000 tonnes of cork a year, with a value of €1.5 billion and employing 30,000 people. Wine corks represent 15% of cork usage by weight but 66% of revenues.
Cork oaks are sometimes planted as individual trees, providing a minor income to their owners. The tree is also sometimes cultivated for ornament. Hybrids with Turkey oak (Quercus cerris) are not uncommon, both where their ranges overlap in the wild in southwest Europe and in cultivation; the hybrid Quercus × hispanica is known as Lucombe oak, for William Lucombe, who first identified it.
Some cork is also produced in eastern Asia from the related Chinese cork oak (Quercus variabilis).
Culture
The cork oak is featured in the city arms of several cities in Portugal, such as the city of Reguengos de Monsaraz, which shows a freshly harvested cork tree.
In 2007, a 2 euro commemorative coin with the motif of a cork oak was issued in Portugal in memory of the Portuguese Presidency of the European Union.
Notable trees
In the Portuguese town of Águas de Moura lies the Sobreiro Monumental (Monumental Cork Oak), also known as 'The Whistler Tree', a tree 236 years old (planted in 1783/1784), over tall and with a trunk that requires at least three people to embrace it. It has been considered a National Monument since 1988, and Guinness World Records lists it as the largest cork tree in the world.
| Biology and health sciences | Fagales | Plants |
344974 | https://en.wikipedia.org/wiki/Chemosynthesis | Chemosynthesis | In biochemistry, chemosynthesis is the biological conversion of one or more carbon-containing molecules (usually carbon dioxide or methane) and nutrients into organic matter using the oxidation of inorganic compounds (e.g., hydrogen gas, hydrogen sulfide) or ferrous ions as a source of energy, rather than sunlight, as in photosynthesis. Chemoautotrophs, organisms that obtain carbon from carbon dioxide through chemosynthesis, are phylogenetically diverse. Groups that include conspicuous or biogeochemically important taxa include the sulfur-oxidizing Gammaproteobacteria, the Campylobacterota, the Aquificota, the methanogenic archaea, and the neutrophilic iron-oxidizing bacteria.
Many microorganisms in dark regions of the oceans use chemosynthesis to produce biomass from single-carbon molecules. Two categories can be distinguished. In the rare sites where hydrogen molecules (H2) are available, the energy available from the reaction between CO2 and H2 (leading to production of methane, CH4) can be large enough to drive the production of biomass. Alternatively, in most oceanic environments, energy for chemosynthesis derives from reactions in which substances such as hydrogen sulfide or ammonia are oxidized. This may occur with or without the presence of oxygen.
Many chemosynthetic microorganisms are consumed by other organisms in the ocean, and symbiotic associations between chemosynthesizers and respiring heterotrophs are quite common. Large populations of animals can be supported by chemosynthetic secondary production at hydrothermal vents, methane clathrates, cold seeps, whale falls, and isolated cave water.
It has been hypothesized that anaerobic chemosynthesis may support life below the surface of Mars, Jupiter's moon Europa, and other planets. Chemosynthesis may have also been the first type of metabolism that evolved on Earth, leading the way for cellular respiration and photosynthesis to develop later.
Hydrogen sulfide chemosynthesis process
Giant tube worms use bacteria in their trophosome to fix carbon dioxide (using hydrogen sulfide as their energy source) and produce sugars and amino acids.
Some reactions produce sulfur:
hydrogen sulfide chemosynthesis:
18H2S + 6CO2 + 3O2 → C6H12O6 (carbohydrate) + 12H2O + 18S
Instead of releasing oxygen gas while fixing carbon dioxide as in photosynthesis, hydrogen sulfide chemosynthesis produces solid globules of sulfur in the process. In bacteria capable of chemoautotrophy (a form a chemosynthesis), such as purple sulfur bacteria, yellow globules of sulfur are present and visible in the cytoplasm.
Discovery
In 1890, Sergei Winogradsky proposed a novel type of life process called "anorgoxydant". His discovery suggested that some microbes could live solely on inorganic matter and emerged during his physiological research in the 1880s in Strasbourg and Zürich on sulfur, iron, and nitrogen bacteria.
In 1897, Wilhelm Pfeffer coined the term "chemosynthesis" for the energy production by oxidation of inorganic substances, in association with autotrophic carbon dioxide assimilation—what would be named today as chemolithoautotrophy. Later, the term would be expanded to include also chemoorganoautotrophs, which are organisms that use organic energy substrates in order to assimilate carbon dioxide. Thus, chemosynthesis can be seen as a synonym of chemoautotrophy.
The term "chemotrophy", less restrictive, was introduced in the 1940s by André Lwoff for the production of energy by the oxidation of electron donors, organic or not, associated with auto- or heterotrophy.
Hydrothermal vents
Winogradsky's suggestion was confirmed nearly 90 years later, when hydrothermal ocean vents were discovered in the 1970s. The hot springs and strange creatures were discovered by Alvin, the world's first deep-sea submersible, in 1977 at the Galapagos Rift. At about the same time, then-graduate student Colleen Cavanaugh proposed chemosynthetic bacteria that oxidize sulfides or elemental sulfur as a mechanism by which tube worms could survive near hydrothermal vents. Cavanaugh later managed to confirm that this was indeed the method by which the worms could thrive, and is generally credited with the discovery of chemosynthesis.
A 2004 television series hosted by Bill Nye named chemosynthesis as one of the 100 greatest scientific discoveries of all time.
Oceanic crust
In 2013, researchers reported their discovery of bacteria living in the rock of the oceanic crust below the thick layers of sediment, and apart from the hydrothermal vents that form along the edges of the tectonic plates. Preliminary findings are that these bacteria subsist on the hydrogen produced by chemical reduction of olivine by seawater circulating in the small veins that permeate the basalt that comprises oceanic crust. The bacteria synthesize methane by combining hydrogen and carbon dioxide.
Chemosynthesis as an innovative area for continuing research
Despite the fact that the process of chemosynthesis has been known for more than a hundred years, its significance and importance are still relevant today in the transformation of chemical elements in biogeochemical cycles. Today, the vital processes of nitrifying bacteria, which lead to the oxidation of ammonia to nitric acid, require scientific substantiation and additional research. The ability of bacteria to convert inorganic substances into organic ones suggests that chemosynthetics can accumulate valuable resources for human needs.
Chemosynthetic communities in different environments are important biological systems in terms of their ecology, evolution and biogeography, as well as their potential as indicators of the availability of permanent hydrocarbon- based energy sources. In the process of chemosynthesis, bacteria produce organic matter where photosynthesis is impossible. Isolation of thermophilic sulfate-reducing bacteria Thermodesulfovibrio yellowstonii and other types of chemosynthetics provides prospects for further research. Thus, the importance of chemosynthesis remains relevant for use in innovative technologies, conservation of ecosystems, human life in general. Sergey Winogradsky helped discover the phenomenon of chemosynthesis.
| Biology and health sciences | Basics | Biology |
345035 | https://en.wikipedia.org/wiki/Halo%20%28optical%20phenomenon%29 | Halo (optical phenomenon) | A halo () is an optical phenomenon produced by light (typically from the Sun or Moon) interacting with ice crystals suspended in the atmosphere. Halos can have many forms, ranging from colored or white rings to arcs and spots in the sky. Many of these appear near the Sun or Moon, but others occur elsewhere or even in the opposite part of the sky. Among the best known halo types are the circular halo (properly called the 22° halo), light pillars, and sun dogs, but many others occur; some are fairly common while others are extremely rare.
The ice crystals responsible for halos are typically suspended in cirrus or cirrostratus clouds in the upper troposphere (), but in cold weather they can also float near the ground, in which case they are referred to as diamond dust. The particular shape and orientation of the crystals are responsible for the type of halo observed. Light is reflected and refracted by the ice crystals and may split into colors because of dispersion. The crystals behave like prisms and mirrors, refracting and reflecting light between their faces, sending shafts of light in particular directions.
Atmospheric optical phenomena like halos were part of weather lore, which was an empirical means of weather forecasting before meteorology was developed. They often do indicate that rain will fall within the next 24 hours, since the cirrostratus clouds that cause them can signify an approaching frontal system.
Other common types of optical phenomena involving water droplets rather than ice crystals include the glory and the rainbow.
History
While Aristotle had mentioned halos and parhelia, in antiquity, the first European descriptions of complex displays were those of Christoph Scheiner in Rome (), Johannes Hevelius in Danzig (1661), and Tobias Lowitz in St Petersburg ().
Chinese observers had recorded these for centuries, the first reference being a section of the "Official History of the Chin Dynasty" (Chin Shu) in 637, on the "Ten Haloes", giving technical terms for 26 solar halo phenomena.
Vädersolstavlan
While mostly known and often quoted for being the oldest color depiction of the city of Stockholm, Vädersolstavlan (Swedish; "The Sundog Painting", literally "The Weather Sun Painting") is arguably also one of the oldest known depictions of a halo display, including a pair of sun dogs. For two hours in the morning of 20 April 1535, the skies over the city were filled with white circles and arcs crossing the sky, while additional suns (i.e., sun dogs) appeared around the Sun.
Light pillar
A light pillar, or sun pillar, appears as a vertical pillar or column of light rising from the Sun near sunset or sunrise, though it can appear below the Sun, particularly if the observer is at a high elevation or altitude. Hexagonal plate- and column-shaped ice crystals cause the phenomenon. Plate crystals generally cause pillars only when the Sun is within 6 degrees of the horizon; column crystals can cause a pillar when the Sun is as high as 20 degrees above the horizon. The crystals tend to orient themselves near-horizontally as they fall or float through the air, and the width and visibility of a sun pillar depend on crystal alignment.
Light pillars can also form around the Moon, and around street lights or other bright lights. Pillars forming from ground-based light sources may appear much taller than those associated with the Sun or Moon. Since the observer is closer to the light source, crystal orientation matters less in the formation of these pillars.
Circular halo
Among the best-known halos is the 22° halo, often just called "halo", which appears as a large ring around the Sun or Moon with a radius of about 22° (roughly the width of an outstretched hand at arm's length). The ice crystals that cause the 22° halo are oriented semi-randomly in the atmosphere, in contrast to the horizontal orientation required for some other halos such as sun dogs and light pillars. As a result of the optical properties of the ice crystals involved, no light is reflected towards the inside of the ring, leaving the sky noticeably darker than the sky around it, and giving it the impression of a "hole in the sky". The 22° halo is not to be confused with the corona, which is a different optical phenomenon caused by water droplets rather than ice crystals, and which has the appearance of a multicolored disk rather than a ring.
Other halos can form at 46° to the Sun, or at the horizon, or around the zenith, and can appear as full halos or incomplete arcs.
Bottlinger's ring
A Bottlinger's ring is a rare type of halo that is elliptical instead of circular. It has a small diameter, which makes it very difficult to see in the Sun's glare and more likely to be noticed around the dimmer subsun, often seen from mountain tops or airplanes. Bottlinger's rings are not well understood yet. It is suggested that they are formed by very flat pyramidal ice crystals with faces at uncommonly low angles, suspended horizontally in the atmosphere. These precise and physically problematic requirements would explain why the halo is very rare.
Other names
In the Cornish dialect of English, a halo around the sun or the moon is called a cock's eye and is an omen of bad weather. The term is related to the Breton word kog-heol (sun cock) which has the same meaning. In Nepal, the halo round the sun is called Indrasabha with a connotation of the assembly court of Lord Indra – the Hindu god of lightning, thunder, and rain.
Artificial halos
The natural phenomena may be reproduced artificially by several means. Firstly, by computer simulations, or secondly by experimental means. Regarding the latter, this occurs when a single crystal is rotated around the appropriate axis/axes, or a chemical approach. A still further and more indirect experimental approach is to find analogous refraction geometries.
Analogous refraction approach
This approach employs the fact that in some cases the average geometry of refraction through an ice crystal may be imitated / mimicked via the refraction through another geometrical object. In this way, the circumzenithal arc, the circumhorizontal arc, and the suncave Parry arcs may be recreated by refraction through rotationally symmetric (i.e. non-prismatic) static bodies. A particularly simple table-top experiment reproduces artificially the colorful circumzenithal and circumhorizontal arcs using a water glass only. The refraction through the cylinder of water turns out to be (almost) identical to the rotationally averaged refraction through an upright hexagonal ice crystal / plate-oriented crystals, thereby creating vividly colored circumzenithal and the circumhorizontal arcs. In fact, the water glass experiment is often confused as representing a rainbow and has been around at least since 1920.
Following Huygens' idea of the (false) mechanism of the 22° parhelia, one may also illuminate (from the side) a water-filled cylindrical glass with an inner central obstruction of half the glasses' diameter to achieve upon projection on a screen an appearance which closely resembles parhelia (cf. footnote [39] in Ref., or see here), an inner red edge transitioning into a white band at larger angles on both sides of the direct transmission direction. However, while the visual match is close, this particular experiment does not involve a fake caustic mechanism and is thus no real analogue.
Chemical approaches
The earliest chemical recipes to generate artificial halos has been put forward by Brewster and studied further by A. Cornu in 1889. The idea was to generate crystals by precipitation of a salt solution. The innumerable small crystals hereby generated will then, upon illumination with light, cause halos corresponding to the particular crystal geometry and the orientation / alignment. Several recipes exist and continue to be discovered. Rings are a common outcome of such experiments. But also Parry arcs have been artificially produced in this way.
Mechanical approaches
Single axis
The earliest experimental studies on halo phenomena have been attributed to Auguste Bravais in 1847. Bravais used an equilateral glass prism which he spun around its vertical axis. When illuminated by parallel white light, this produced an artificial parhelic circle and many of the embedded parhelia. Similarly, A. Wegener used hexagonal rotating crystals to produce artificial subparhelia. In a more recent version of this experiment, many more embedded parhelia have been found using commercially available hexagonal BK7 glass crystals. Simple experiments like these can be used for educational purposes and demonstration experiments. Unfortunately, using glass crystals one cannot reproduce the circumzenithal arc or the circumhorizontal arc due to total internal reflections preventing the required ray-paths when .
Even earlier than Bravais, the Italian scientist F. Venturi experimented with pointed water-filled prisms to demonstrate the circumzenithal arc. However, this explanation was replaced later by the CZA's correct explanation by Bravais.
Artificial ice crystals have been employed to create halos which are otherwise unattainable in the mechanical approach via the use of glass crystals, e.g. circumzenithal and circumhorizontal arcs. The use of ice crystals ensures that the generated halos have the same angular coordinates as the natural phenomena. Other crystals such as sodium fluoride (NaF) also have a refractive index close to ice and have been used in the past.
Two axes
In order to produce artificial halos such as the tangent arcs or the circumscribed halo one should rotate a single columnar hexagonal crystal about 2 axes. Similarly, the Lowitz arcs can be created by rotating a single plate crystal about two axes. This can be done by engineered halo machines. The first such machine was constructed in 2003; several more followed. Putting such machines inside spherical projection screens, and by the principle of the so-called sky transform, the analogy is nearly perfect. A realization using micro-versions of the aforementioned machines produces authentic distortion-free projections of such complex artificial halos. Finally, superposition of several images and projections produced by such halo machines may be combined to create a single image. The resulting superposition image is then a representation of complex natural halo displays containing many different orientation sets of ice prisms.
Three axes
The experimental reproduction of circular halos is the most difficult using a single crystal only, while it is the simplest and typically achieved one using chemical recipes. Using a single crystal, one needs to realize all possible 3D orientations of the crystal. This has recently been achieved by two approaches. The first one using pneumatics and a sophisticated rigging, and a second one using an Arduino-based random walk machine which stochastically reorients a crystal embedded in a transparent thin-walled sphere.
Gallery
| Physical sciences | Atmospheric optics | Earth science |
345286 | https://en.wikipedia.org/wiki/Gel%20permeation%20chromatography | Gel permeation chromatography | Gel permeation chromatography (GPC) is a type of size-exclusion chromatography (SEC), that separates high molecular weight or colloidal analytes on the basis of size or diameter, typically in organic solvents. The technique is often used for the analysis of polymers. As a technique, SEC was first developed in 1955 by Lathe and Ruthven. The term gel permeation chromatography can be traced back to J.C. Moore of the Dow Chemical Company who investigated the technique in 1964. The proprietary column technology was licensed to Waters Corporation, who subsequently commercialized this technology in 1964. GPC systems and consumables are now also available from a number of manufacturers. It is often necessary to separate polymers, both to analyze them as well as to purify the desired product.
When characterizing polymers, it is important to consider their size distribution and dispersity (Đ) as well their molecular weight. Polymers can be characterized by a variety of definitions for molecular weight including the number average molecular weight (Mn), the weight average molecular weight (Mw) (see molar mass distribution), the size average molecular weight (Mz), or the viscosity molecular weight (Mv). GPC allows for the determination of Đ as well as Mv and, based on other data, the Mn, Mw, and Mz can be determined.
How it works
GPC is a type of chromatography in which analytes are separated, based on their size or hydrodynamic volume (radius of gyration). This differs from other chromatographic techniques, which depend upon chemical or physical interactions between the mobile and stationary phases to separate analytes. Separation occurs via the use of porous gel beads packed inside a column (see stationary phase (chemistry)). The principle of separation relies on the differential exclusion or inclusion of the macromolecules by the porous gel stationary phase. Larger molecules are excluded from entering the pores and elute earlier, while smaller molecules can enter the pores, thus staying longer inside the column. The entire process takes place without any interaction of the analytes with the surface of the stationary phase. The smaller analytes relative to the pore sizes can permeate these pores and spend more time inside the gel particles, increasing their retention time. Conversely, larger analytes relative to the pores sizes spend little if any time inside the column, hence they elute sooner. Each type of column has a range of molecular weights that can be separated, according to their pores sizes. If an analyte is too large relative to the column's pores, it will not be retained at all and will be totally excluded; conversely, if the analyte is small relative to the pores sizes, it will be totally permeating. Analytes that are totally excluded, elute with the free volume outside around the particles (Vo), the total exclusion limit, while analytes that are completely delayed, elute with the solvent, marking the total permeation volume of the column, including also the solvent held inside the pores (Vi). The total volume can be considered by the following equation, where Vg is the volume of the polymer gel and Vt is the total volume:
As can be inferred, there is a limited range of molecular weights that can be separated by each column, therefore the size of the pores for the packing should be chosen according to the range of molecular weight of analytes to be separated. For polymer separations the pore sizes should be on the order of the polymers being analyzed. If a sample has a broad molecular weight range it may be necessary to use several GPC columns with varying pores volumes in tandem to resolve the sample fully.
Application
GPC is often used to determine the relative molecular weight of polymer samples as well as the distribution of molecular weights. What GPC truly measures is the molecular volume and shape function as defined by the intrinsic viscosity. If comparable standards are used, this relative data can be used to determine molecular weights within ± 5% accuracy. Polystyrene standards with dispersities of less than 1.2 are typically used to calibrate the GPC. Unfortunately, polystyrene tends to be a very linear polymer and therefore as a standard it is only useful to compare it to other polymers that are known to be linear and of relatively the same size.
Material and methods
Instrumentation
Gel permeation chromatography is conducted almost exclusively in chromatography systems. The experimental design is not much different from other techniques of High Performance liquid chromatography. Samples are dissolved in an appropriate solvent, in the case of GPC these tend to be organic solvents and after filtering the solution it is injected onto a column. The separation of multi-component mixture takes place in the column. The constant supply of fresh eluent to the column is accomplished by the use of a pump. Since most analytes are not visible to the naked eye a detector is needed. Often multiple detectors are used to gain additional information about the polymer sample. The availability of a detector makes the fractionation convenient and accurate.
Gel
Gels are used as stationary phase for GPC. The pore size of a gel must be carefully controlled in order to be able to apply the gel to a given separation. Other desirable properties of the gel forming agent are the absence of ionizing groups and, in a given solvent, low affinity for the substances to be separated. Commercial gels like PLgel & Styragel (cross-linked polystyrene-divinylbenzene), LH-20 (hydroxypropylated Sephadex), Bio-Gel (cross-linked polyacrylamide), HW-20 & HW-40 (hydroxylated methacrylic polymer), and agarose gel are often used based on different separation requirements.
Column
The column used for GPC is filled with a microporous packing material. The column is filled with the gel. Since the total penetration volume is the maximum volume permeated by the analytes, and there is no retention on the surface of the stationary phase, the total column volume is usually large, relatively to the sample volume.
Eluent
The eluent (mobile phase) should be the appropriate solvent to dissolve the polymer, should not interfere with the response of the polymer analyzed, and should wet the packing surface and make it inert to interactions with the polymers. The most common eluents for polymers that dissolve at room temperature GPC are tetrahydrofuran (THF), o-dichlorobenzene and trichlorobenzene at 130–150 °C for crystalline polyalkynes and hexafluoroisopropanol (HFIP) for crystalline condensation polymers such as polyamides and polyesters.
Pump
There are two types of pumps available for uniform delivery of relatively small liquid volumes for GPC: piston or peristaltic pumps. The delivery of a constant flow free of fluctuations is especially important to the precision of the GPC analysis, as the flow-rate is used for the calibration of the molecular weight, or diameter.
Detector
In GPC, the concentration by weight of polymer in the eluting solvent may be monitored continuously with a detector. There are many detector types available and they can be divided into two main categories. The first is concentration sensitive detectors which includes UV-VIS absorption, differential refractometer (DRI) or refractive index (RI) detectors, infrared (IR) absorption and density detectors. The second category is molecular weight sensitive detectors, which include low angle light scattering detectors (LALLS) and multi angle light scattering (MALS). The resulting chromatogram is therefore a weight distribution of the polymer as a function of retention volume. The most sensitive detector is the differential UV photometer and the most common detector is the differential refractometer (DRI). When characterizing copolymer, it is necessary to have two detectors in series. For accurate determinations of copolymer composition at least two of those detectors should be concentration detectors. The determination of most copolymer compositions is done using UV and RI detectors, although other combinations can be used.
Data analysis
Gel permeation chromatography (GPC) has become the most widely used technique for analyzing polymer samples in order to determine their molecular weights and weight distributions. Examples of GPC chromatograms of polystyrene samples with their molecular weights and dispersities are shown on the left.
Benoit and co-workers proposed that the hydrodynamic volume, Vη, which is proportional to the product of [η] and M, where [η] is the intrinsic viscosity of the polymer in the SEC eluent, may be used as the universal calibration parameter. If the Mark–Houwink–Sakurada constants K and α are known (see Mark–Houwink equation), a plot of log [η]M versus elution volume (or elution time) for a particular solvent, column and instrument provides a universal calibration curve which can be used for any polymer in that solvent. By determining the retention volumes (or times) of monodisperse polymer standards (e.g. solutions of monodispersed polystyrene in THF), a calibration curve can be obtained by plotting the logarithm of the molecular weight versus the retention time or volume. Once the calibration curve is obtained, the gel permeation chromatogram of any other polymer can be obtained in the same solvent and the molecular weights (usually Mn and Mw) and the complete molecular weight distribution for the polymer can be determined. A typical calibration curve is shown to the right and the molecular weight from an unknown sample can be obtained from the calibration curve.
Advantages
As a separation technique, GPC has many advantages. First of all, it has a well-defined separation time due to the fact that there is a final elution volume for all unretained analytes. Additionally, GPC can provide narrow bands, although this aspect of GPC is more difficult for polymer samples that have broad ranges of molecular weights present. Finally, since the analytes do not interact chemically or physically with the column, there is a lower chance for analyte loss to occur. For investigating the properties of polymer samples in particular, GPC can be very advantageous. GPC provides a more convenient method of determining the molecular weights of polymers. In fact most samples can be thoroughly analyzed in an hour or less. Other methods used in the past were fractional extraction and fractional precipitation. As these processes were quite labor-intensive molecular weights and mass distributions typically were not analyzed. Therefore, GPC has allowed for the quick and relatively easy estimation of molecular weights and distribution for polymer samples
Disadvantages
There are disadvantages to GPC, however. First, there is a limited number of peaks that can be resolved within the short time scale of the GPC run. Also, as a technique GPC requires around at least a 10% difference in molecular weight for a reasonable resolution of peaks to occur. In regards to polymers, the molecular masses of most of the chains will be too close for the GPC separation to show anything more than broad peaks. Another disadvantage of GPC for polymers is that filtrations must be performed before using the instrument to prevent dust and other particulates from ruining the columns and interfering with the detectors. Although useful for protecting the instrument, there is the possibility of the pre-filtration of the sample removing higher molecular weight sample before it can be loaded on the column. Another possibility to overcome these issues is the separation by field-flow fractionation (FFF).
Orthogonal methods
Field-flow fractionation (FFF) can be considered as an alternative to GPC, especially when particles or high molar mass polymers cause clogging of the column, shear degradation is an issue or agglomeration takes place but cannot be made visible. FFF is separation in an open flow channel without having a static phase involved so no interactions occur. With one field-flow fractionation version, thermal field-flow fractionation, separation of polymers having the same size but different chemical compositions is possible.
| Physical sciences | Chromatography | Chemistry |
345330 | https://en.wikipedia.org/wiki/Iridaceae | Iridaceae | Iridaceae () is a family of plants in order Asparagales, taking its name from the irises. It has a nearly global distribution, with 69 accepted genera with a total of c. 2500 species. It includes a number of economically important cultivated plants, such as species of Freesia, Gladiolus, and Crocus, as well as the crop saffron.
Members of this family are perennial plants, with a bulb, corm or rhizome. The plants grow erect, and have leaves that are generally grass-like, with a sharp central fold. Some examples of members of this family are the blue flag and yellow flag.
Etymology
The family name comes from the genus Iris, the family's largest and best-known genus in Europe. This genus dates from 1753, when it was coined by Swedish botanist, Carl Linnaeus. Its name derives from the Greek goddess, Iris, who carried messages from Olympus to earth along a rainbow, whose colors were seen by Linnaeus in the multi-hued petals of many of the species.
Taxonomy
Iridaceae is currently recognized as nested in the Asparagales order but was traditionally grouped with Liliales. Iridaceae was previously divided into four subfamilies but results from phylogenetic analysis suggested an additional three could be recognized. These differences in circumscription are a result of homoplastic traits, including asymmetric corms, woody corm covering, exclusion of the vascular trace during ovule development, and leaf margin. Molecular clock analyses have supported initial cladogenesis in Antarctica-Australasia 82 mya from a Doryanthaceae ancestor. The distribution of subfamilies in Iridaceae is considered to be phylogenetically structured, with all neotropical species belonging to one subfamily, the Irdoideae.
Crocoideae
Subfamily Crocoideae is one of the major subfamilies in the family Iridaceae. It contains many genera, including Afrocrocus, Babiana, Chasmanthe, Crocosmia, Crocus, Cyanixia, Devia, Dierama, Duthiastrum, Freesia, Geissorhiza, Gladiolus, Hesperantha, Ixia, Lapeirousia, Melasphaerula, Micranthus, Pillansia, Romulea, Sparaxis, Savannosiphon, Syringodea, Thereianthus, Tritonia, Tritoniopsis, Xenoscapa and Watsonia.
They are mainly from Africa, but includes members from Europe and Asia. The rootstock is usually a corm, they have blooms which sometimes have scent, are collected in inflorescence and contain six tepals. The nectar is produced mostly in the base of the bloom from the glands of the ovary, which is where the flower forms a tube-like end. In some species there is no such end and the plant only provides pollen to pollinating insects. Members of this subfamily have the sword-shaped leaves typical of Iridaceae.
Isophysidoideae
Subfamily Isophysidoideae is monotypic, only containing Isophysis from Tasmania. It is the only member of the family with a superior ovary, and it grows a solitary star-like, yellow to brownish flower. It is also sister to all other extant taxa of Iridaceae, diverging 66mya.
Nivenioideae and allies
Subfamily Nivenioideae contained six genera from South Africa, Australia and Madagascar, including the core genera and only true shrubs in the family (Klattia, Nivenia and Witsenia). Upon phylogenetic analysis, subfamily Crocoideae is always found nested within Nivenioideae, leading to it not being a monophyletic taxon. A revised description of these groups led to the description of Aristea, Geosiris, and Patersonia each as separate subfamilies, retaining a core, monophyletic Nivenioideae. It is now distinguished as being evergreen shrubs with monocot-type secondary thickening, shield shaped seeds, and paired rhipidia with only one to two flowers in each cluster.
Iridoideae
Subfamily Iridoideae has the widest geographic distribution and is divided into four tribes and one sister genus: Irideae, Sisyrichieae, Trimezieae, Tigridieae,and Diplarreneae. Iridoideae is differentiated from the other subfamilies by having very short-lived flowers, nectaries on the perianth, and long branching styles. Excluding the Irideae, the evolution of oil-producing trichomes, called elaiophores, have been gained and lost in each of the tribes attracting oil bees. The genus Diplarreneae is sister to the rest of the subfamily and is unique to Iridoideae in having zygomorphic flowers and stamens with unequal height. Irideae represents the Old World portion of the subfamily but include several genera that diversified in North America, such as Iris. They are distinguishable with the presence of flattened anthers pressed to the style, petaloid crests, and schlerenchyma tissue along the margins of leaves. Sisyrichieae is noted for having long style branches that may interlace with stamens, partially fused filaments, and the lack of oxaloacetate crystals in leaves. Trimezieae is the smallest tribe with two to four genera, noted for the presence of large rhizomes or corms rather than bulbs as well as a thickened midrib. Several species with ornamented or iris-like flowers also possess a specialized method of forcing pollen onto heavy pollinators with hinged petals. Tigridieae are distinguished for their large bulbous rootstock and plicate, decidious leaves. The number of genera and whether any morphology can distinguish between them has been debated.
Ecology
Members of Iridaceae occur in a great variety of habitats. Gladiolus gueinzii occurs on the seashore just above the high tide mark within reach of the spray. Most species are adapted to seasonal climates that have a pronounced dry or cold period unfavorable for plant growth and during which the plants are dormant. As a result, most species are deciduous. Evergreen species are restricted to subtropical forests or savanna, temperate grasslands and perennially moist fynbos. A few species grow in marshes or along streams and some even grow only in the spray of seasonal waterfalls.
Members of the subfamilies Crocoideae and Nivenioideae first began cladogenesis in arid conditions in Africa, accelerating for Crocoideae as the Mediterranean climate emerged in Southern Africa. A similar process occurred for the tribe Tigridieae in Iridoideae following long-distance dispersal from South to North America, resulting in high levels of endemism. In the tribe Sisyrichieae, the continued formation of the Andes supported the movement to lower elevations along the Atlantic.
The aerial portions of deciduous species die back when the bulb or corm enters dormancy. The plants thus survive periods that are unfavorable for growth by retreating underground. This is particularly useful in grasslands and fynbos, which are adapted to regular burning in the dry season. At this time the plants are dormant and their bulbs or corms are able to survive the heat of the fires underground. Veld fires clear the soil surface of competing vegetation, as well as fertilize it with ash. With the arrival of the first rains, the dormant corms are ready to burst into growth, sending up flowers and stems before they can be shaded out by other vegetation. Many grassland and fynbos irids flower best after fires and some fynbos species will only flower in the season after a fire.
The majority of Iridaceae are pollinated by Hymenoptera, frequently by single species or a small group of species. These tight relationships found in individual species of Iridaceae, especially in Gladiolus, were the inspiration for the description of pollinator syndromes. Pollinators include various species of solitary bees, as well as sunbirds, long-proboscid flies (such as Moegistorhynchus longirostris), butterflies, and night moths. Ancestrally, flowers were zygomorphic, as in Crocoideae, with contrasting nectary locations for pollinators. Flowers may present nectar and pollen rewards to visitors, but some genera may only offer nectar such as in Gladious and Watsonia. Species of Ferraria produce putrid smells, floral cups, and dark mottled perianth in order to attract Diptera. Members of Iridoideae and Nivenioideae have radially symmetric trumpet-like flowers that secrete large amounts of nectar. This novel morphology enabled additional floral complexity and rapid evolution of pollinator relationships, as frequently as a new relationship over 5 speciations. New World Iridoideae represent one of the largest clades offering oil to pollinators, ranging from forced pollination using hinged petals to frequent failure to pollinate. Most of the variability in flowers occurs between subfamilies, including infloresence structure, i.e. rhipidia, panicle, or spike, and floral longevity, i.e. less than one day to five days. Some members of the tribe Irideae have flowers functioning as meranthia, or developing as three separate zygomorphic units that pollinators visit individually.
List of genera
69 genera have been recognized in the family, with a total of 2597 species described. The Afrotropical realm, and in particular South Africa, have the greatest diversity of genera.
| Biology and health sciences | Monocots | null |
345356 | https://en.wikipedia.org/wiki/Skate%20%28fish%29 | Skate (fish) | Skates are cartilaginous fish belonging to the family Rajidae in the superorder Batoidea of rays. More than 150 species have been described, in 17 genera. Softnose skates and pygmy skates were previously treated as subfamilies of Rajidae (Arhynchobatinae and Gurgesiellinae), but are now considered as distinct families. Alternatively, the name "skate" is used to refer to the entire order of Rajiformes (families Anacanthobatidae, Arhynchobatidae, Gurgesiellidae and Rajidae).
Members of Rajidae are distinguished by a stiff snout and a rostrum that is not reduced.
Taxonomy and systematics
Evolution
Skates belong to the ancient lineage of cartilaginous fishes. Fossil denticles (tooth-like scales in the skin) resembling those of today's chondrichthyans date at least as far back as the Ordovician, with the oldest unambiguous fossils of cartilaginous fish dating from the middle Devonian. A clade within this diverse family, the Neoselachii, emerged by the Triassic, with the best-understood neoselachian fossils dating from the Jurassic. This clade is represented today by sharks, sawfish, rays and skates.
The body plan of skates is caused by skate-specific genomic rearrangements that have altered the three-dimensional regulatory landscape of genes. These changes arose about 286–221 million years ago when skates diverged from sharks.
Classification
The skate belongs to the class Chondrichthyes. This class consists of all the cartilaginous fishes, including sharks and stingrays. Chondrichthyes is divided into two subclasses; of which Elasmobranchii includes skates, rays, and sharks. Skates are the most diverse elasmobranch group, comprising over 20% of the known species. The number of species is likely to increase as taxonomic issues are resolved and new species are identified. Thefollowing genera recognized in the family Rajidae:
Skates have more valid species than any other group of cartilaginous fishes. Since 1950, 126 new species of skates have been discovered. Five scientists take credit for the rapid increase of findings. The Rajidae are considered monophyletic because of their similarity in appearance. There are 18 genera and about 250 valid species. However, there is little information about the diets of about 24% of these species. There are at least 45 dubious species of skates worldwide.
Description
General Batoidea characteristics
Skates are cartilaginous fishes like other Chondrichthyes, however, skates, like rays and other Rajiformes, have a flat body shape with flat pectoral fins that extend the length of their body. This structure creates power for forward propulsion, providing the emergence swimming capabilities that enabled skates to colonize the sea floor.
A large portion of the skate's dorsal body is covered by rough skin made of placoid scales. Placoid scales have a pointed tip that is oriented caudally and are homologous to teeth. Their mouths are located on the underside of the body, with a jaw suspension common to Batoids known as euhyostyly. Skate's gill slits are located ventrally as well, but dorsal spiracles allow the skate to be partially buried in floor sediment and still complete respiratory exchange. Also located on the dorsal side of the skate are their two eyes which allow for predator awareness. In addition to their pectoral fins, skates have a first and second dorsal fin, caudal fin and paired pelvic fins. Distinct from their rhomboidal shape is a long fleshy slender tail. While skate anatomy is similar to other Batoidea, features such as their electric organ and mermaid's purse create clear distinctions.
Skate specific characteristics
Mermaid's purse
Skates produce their young in an egg case called a mermaid's purse. These egg cases have distinct characteristics that are individualized to each species. This makes a great tool for identifying different species of skates. One of these identifiable structures is the keel. The keel is a flexible ridge that runs along the outside of the structure. Another characteristic is the number of embryos in the egg case. Some species contain only one embryo while others can have up to seven. The size of the fibrous shell around the case is another characteristic. Some species have thick layers on the exterior.
Electric organ
The electric organ is a characteristic exclusive to aquatic species. Among the Chondrichthyes, the only groups to possess electric organs are the electric ray and the skates. Unlike many other electrogenic fishes, skates are unique in having paired electric organs which run longitudinally through the tail in the lateral musculature of the notochord. The impulses put out by the electric organs of the skate are considered to be weak, asynchronous, long-lasting signals. Although the anatomy of the skate's electric organ is well described, its function is poorly understood. Some research suggests the electric impulses are too weak to be a mechanism used for defense or hunting. It is also too irregular to be useful for electrolocation purposes. The most reasonable explanation in the literature suggests that the electric organ discharges may be used as a form of communication used for reproduction purposes.
Distribution and habitats
Skates are primarily found from the intertidal down to depths greater than . They are most commonly found along outer continental shelves and upper slopes. They are typically more diverse at higher latitudes and in deep-water. In fact, skates are the only cartilaginous fish taxon to exhibit more diversity of species at higher latitudes. A cool, temperate to polar water in the deep sea can be a favorable environment for skates. As the water becomes more shallow and warmer, skates are seen to be replaced by stingrays. Skates are absent from brackish and freshwater environments. However, there is a single estuarine species that has been found in Tasmania, Australia. Also, the Connecticut Department of Environmental Protection has caught and studied skates within the Long Island Sound estuary. Some skate fauna have been found inhabiting areas of rock cobble and high rocky relief.
Behavior and ecology
Reproduction
Skates mate at the same nursery ground each year. In order to fertilize the egg, males use claspers, a structure attached to the pelvic fins. The claspers allow them to direct the flow of semen into the female's cloaca. Skates are oviparous, meaning they lay eggs with very little development in the mother. This is one major difference from rays, which are viviparous, meaning they give birth to live young. When a female skate is fertilized, a protected case forms around the embryo called an egg case, or more commonly mermaid's purse. This egg case is then deposited out of the mother's body onto the ocean floor where the skates develop for up to 15 months before they enter the external environment.
Diet and feeding
The majority of skates feed on bottom dwelling animals, such as shrimp, crab, oyster, clams, and other invertebrates. To feed on these animals they have grinding plates in their mouths. Skates are an influential part of the food webs of demersal marine communities. They utilize similar resources to those of other upper trophic-level marine predators, such as seabirds, marine mammals, and sharks. The flattened body shape, ventral eyes and well developed spiracles of the skate allows them to live benthically, buried in the sediment or using a longitudinal undulation of the pectoral fins known as Rajiform locomotion to glide along the water floor. Current research suggests that some species of skates, in addition to their Rajiform locomotion, use their pelvic fins to perform ambulatory locomotion. This form of locomotion performed by the skate is being explored as a possible origin for our own development of walking by looking for similar neural pathways used for movement between skates and animals walking on land.
Skates versus stingrays
Skates are like stingrays in that they have five pairs of gill slits that are located ventrally, which means on the underside of their body (unlike sharks that have their gills located on their sides). Skates and rays both have pectoral fins that are flat and expanded, which are typically fused to the head. Both skates and stingrays typically have their eyes on top of their head. Skates also share similar feeding habits with rays.
Skates are different from rays in that they lack a whip-like tail and stinging spines. However, some skates have electric organs located in their tail. The main difference between skates and rays is that skates lay eggs, whereas rays give birth to live young.
Moreover, skates can be more abundant than rays, and are fished for food in some parts of the world.
Conservation
Skates have slow growth rates and, since they mature late, low reproductive rates. As a result, skates are vulnerable to overfishing and appear to have been overfished and are suffering reduced population levels in many parts of the world.
| Biology and health sciences | Batoidea | null |
345366 | https://en.wikipedia.org/wiki/Luna%20moth | Luna moth | The luna moth (Actias luna), also called the American moon moth, is a Nearctic moth in the family Saturniidae, subfamily Saturniinae, a group commonly named the giant silk moths.
The moth has lime-green wings and a white body. Its caterpillars are also green. Its typical wingspan is roughly , but wingspans can exceed , ranking the species as one of the larger moths in North America.
Across Canada, it has one generation per year, with the winged adults appearing in late May or early June, whereas farther south it will have two or even three generations per year, the first appearance as early as March in southern parts of the United States.
As defense mechanisms, larvae emit clicks as a warning and can also regurgitate intestinal contents, confirmed as having a deterrent effect on a variety of predators. The elongated tails of the hindwings are thought to confuse the echolocation detection used by predatory bats.
A parasitoid fly deliberately introduced to North America as a biological pest control for the invasive species spongy moth (also known as gypsy moth) appears to have had a negative impact on luna moths and other native moths.
Description
Eggs, attached in small groups to undersides of leaves, are mottled white and brown, slightly oval, and roughly 1.5 millimeters in diameter. Larvae are primarily green, with sparse hairs. The first instar, emerging from the egg, reaches a length of , the second , the third and the fourth . The fifth (final) instar grows to approximately in length. Small, colorful dots – yellow or magenta – may line the sides of the fourth and fifth instars. The larvae may take on a reddish-brown color just prior to cocooning. Fifth-instar larvae descend to the ground and use silk to bind dead leaves around the cocoon.
The imagoes (winged, sexually mature), often referred to as 'adult moths,' emerge from the pupae with the wings small, crumpled and held close to the body. Over several hours the wings will enlarge to full size. Wingspan is typically , and in rare instances as much as . Females and males are similar in size and appearance: green wings, eyespots on both forewings and hind wings, and long, sometimes somewhat twisted tails extending from the back edge of the hindwings. Bodies are white and hairy. Adults have vestigial mouthparts and do not feed. Energy is from fat stores created while a caterpillar. The forward edge of the forewing is dark-colored and thick, tapering in thickness from the thorax to the wing tip. Its color can range from maroon to brown. The eyespots, one per wing, are oval in shape on the forewings and round on the hindwings. Each eyespot can have arcs of black, blue, red, yellow, green or white. The eyespots are thought to confuse potential predators.
There are some sex-determined and regional differences in appearance. Females will have a larger abdomen compared to males because it contains 200–400 eggs. Both sexes have antennae, but on the male, much longer and wider. Wing color is blue-green in the north and for the overwintering generation in the central and southern states; second and third generation wing color has more of a yellow-green tint.
Etymology
Described and named Phalena plumata caudata by James Petiver in 1700, this was the first North American saturniid to be reported in the insect literature. The initial Latin name, which roughly translates to "brilliant, feather tail", was replaced when Carl Linnaeus described the species in 1758 in the tenth edition of Systema Naturae, and renamed it Phalaena luna, later Actias luna, with luna derived from Luna, the Roman moon goddess. The common name became "Luna moth". Several other North American giant silk moths were also given species names after Roman or Greek mythology.
Distribution
The Luna moth is found in North America, from east of the Great Plains in the United States – Florida to Maine, and from Saskatchewan eastward through central Quebec to Nova Scotia in Canada. Luna moths are also rarely found in Western Europe as vagrants.
Life cycle
Based on the climate in which they live, Luna moths produce different numbers of generations per year. In Canada and northern regions of the United States, they are univoltine, meaning one generation per year. Life stages are approximately 10 days as eggs, 6–7 weeks as larvae, 2–3 weeks as pupae, finishing with one week as winged adults appearing in late May or early June. In the mid-Atlantic states, the species is bivoltine, characterized by two generations per year. In contrast, farther south, they are trivoltine, producing three generations within the same time frame. In the central states, the first generation appears in April, the second in July. Even farther south, the first generation appears as early as March, with the second and third spaced eight to ten weeks later.
Eggs
Females lay 200–400 eggs, singly or in small groups, on the underside of leaves of the tree species preferred by the larvae. Egg laying starts the evening after mating is completed and goes on for several days. Eggs hatch in about a week.
Larvae
Each instar – the period between molts – generally takes about 4–10 days. There are five instars before cocooning. At the end of each instar, a small amount of silk is placed on the major vein of a leaf and the larva undergoes apolysis, then ecdysis (molting), leaving the old exoskeleton behind. Sometimes the shed exoskeleton is eaten. Newly hatched, this caterpillar constantly munches on the leaves of walnut, hickory, sweetgum, and paper birch trees. Each instar is green, though the first two instars do have some variation in which some larvae will have black underlying splotches on their dorsal side. The final instar grows to approximately in length. All five instar stages possess green spines on the dorsal surface. These spines do not sting, but can still cause irritation upon contact. This is a tree-dwelling species. Larvae stay on the same tree where they hatched until it is time to descend to the ground to make a cocoon. When females emerge from cocoons, they fly to preferred tree species, emit pheromones, and wait there for males to find them. Although some larvae in the family Saturniidae are known to be poisonous, those of A. luna are not. The spines, or setae, located on the thoracic and abdominal segments have no chemical component to them.
Pupae
The Luna moth pupates after spinning a silk cocoon, which is thin and single layered. Shortly before pupation, the final, fifth-instar caterpillar will engage in a "gut dump" where any excess water and intestinal contents are expelled. As pupae, this species is more physically active than most moths. When disturbed, the moths will wiggle within their pupal cases, producing a noise. Pupation takes approximately two weeks unless the individual is in diapause over winter, in which case the pupal stage takes about nine months. The mechanisms triggering diapause are generally a mixture of genetic triggers, duration of sunlight and temperature. The pupae have chitinous spurs near the base of the forewings. By vigorously moving about within the cocoon, these spurs tear a circular opening from which the imago emerges, the silk of the cocoon having also been weakened by the secretion of cocoonase, a protein-digesting enzyme.
Imago (winged)
Pupae transition to winged state after receiving external signals in the form of temperature change. When the adult Luna moths emerge from their pupae, their abdomens are swollen and their wings are small, soft and wet. The first few hours of adult life will be spent pumping hemolymph (invertebrates' equivalent to blood) from the abdomen into the wings. The moths must wait for the wings to dry and harden before being able to fly. This process can take 2–3 hours to complete. Luna moths are not rare, but are rarely seen due to their very brief (7–10 day) adult lives and nocturnal flying time. As with all giant silk moths, the adults only have vestigial mouthparts and no digestive system and therefore do not eat in their adult form. Instead, they rely on energy stored during their caterpillar stage. In regions where there are two or three generations per year, the second and third may have wing coloration that is more of a yellow-green compared to the first generation of the year.
Mating
Giant silk moths have in common a mating process wherein the females, at night, release volatile sex pheromones, which the males, flying, detect via their large antennae. Males can detect these molecules at a distance of several miles, and then fly in the direction the wind is coming from until reaching the female. Luna moth females mate with the first males to find them, a process that typically starts after midnight and takes several hours. Researchers extracted three chemical compounds from the pheromone gland of unmated Luna moth females and identified one major and two minor aldehyde compounds designated E6,Z11-18:Ald, E6-18:Ald and Z11-18:Ald. The same compounds were also synthesized. Field experiments with both unmated females and the synthesized compounds confirmed that E6, Z11-18:Ald was the major sex pheromone, attraction augmented by the addition of E6-18:Ald but not by Z11-18:Ald. The authors mentioned that no other moth species were attracted to either the unmated females or the synthesized products, confirming that the pheromone is species-specific, at least for the sites and dates where it was tested.
Gallery of life cycle
Close-up images
Predators and parasites
Some species of giant silk moth larvae are known to make clicking noises when attacked by rubbing their serrated mandibles together. These clicks are audible to humans and extend into ultrasound frequencies audible to predators. Clicks are thought to be a form of aposematic warning signaling, made prior to predator-deterring regurgitation of intestinal contents. Luna moth larvae click and regurgitate, with the regurgitated material confirmed as being a predator deterrent against several species.
Imagos (winged adults) of this and related night-flying Actias species, collectively referred to as "moon moths", have long hindwing tails. A "false target" hypothesis holds that the tails evolved to reduce predation risk by bats which use echolocation to locate prey. The moths use the spinning hindwing tails to fool bats into attacking nonessential appendages, with success occurring over 55% of the time. Experiments were conducted with Luna moths with intact wings and with the tails removed. With intact wings, a majority of the attacking bats contacted the hindwing tails rather than the body of the moth; only 35% of intact moths were caught versus 81% for those with clipped tails. The results of this experiment support echolocation distortion as an effective countermeasure.
The parasitoid tachinid fly Compsilura concinnata native to Europe was deliberately introduced to the United States throughout much of the 20th century as a biological control for the gypsy moth (Lymantria dispar) (also known as the "spongy moth"). Researchers reported that when Luna moth larvae were placed outside for about a week and then collected and returned to the laboratory, four parasitic species emerged, the most common being C. concinnata. The researchers concluded that this parasitoid fly causes collateral damage to Luna moth populations.
Luna moth larvae have displayed defenses against predators in late instars by developing spines once they reach about 3 cm in length. Unlike other species such as Automeris io, which have chemical defenses much earlier in the larval stage, the Luna moth larvae are left largely defenseless until it reaches this length. However, the absence of a chemical defense allows for the shortening of the larval stage. Automeris io has a larval stage at least twice as long on average as Actias luna, leaving it vulnerable to parasitism.
Host plants
The larvae of Luna moths feed on several different species of broadleaf trees. The larvae do not reach population densities sufficient to cause significant damage to their host trees. Tuskes listed white birch (Betula papyrifera), American persimmon (Diospyros virginiana), American sweet gum (Liquidambar styraciflua), plus several species of hickory (Carya), walnut (Juglans) and sumac (Rhus) as host plants for the caterpillars. Other tree species have been identified as suitable for Actias luna larvae, but a feeding experiment that also included black cherry, eastern cottonwood, quaking aspen, white willow, red oak, white oak and tulip tree reported very poor survival on these seven tree species even though older literature had identified them as hosts. The author suggested that host plant utilization may differ regionally, so that larvae collected from one region may not tolerate host plants readily consumed in another region. Biochemical detoxification of host plant defensive chemicals by digestive system enzymes may be a factor in regional host plant specialization. Juglone is a chemical compound common to walnut and hickory which most insects find a deterrent or even toxic. Luna moth larvae have higher concentrations of juglone-neutralizing digestive system enzymes compared to other lepidoptera, and concentrations were even higher when larvae were fed walnut or hickory leaves versus white birch or American sweet gum. This suggests evolutionary and inducible adaptations to allow consumption of certain host plants.
In popular culture
The Luna moth appeared on a first class United States postage stamp issued in June 1987. Although more than two dozen butterflies have been so honored, as of 2019 this is the only moth.
| Biology and health sciences | Lepidoptera | Animals |
345526 | https://en.wikipedia.org/wiki/DEET | DEET | N,N-Diethyl-meta-toluamide, also called diethyltoluamide or DEET (, from DET, the initials of di- + ethyl + toluamide), is the oldest, one of the most effective, and most common active ingredients in commercial insect repellents. It is a slightly yellow oil intended to be applied to the skin or to clothing and provides protection against mosquitoes, flies, ticks, fleas, chiggers, leeches, and many other biting insects.
Effectiveness
DEET is effective against a variety of invertebrates, including ticks, flies, mosquitos, and some parasitic worms.
A 2018 systematic review found no consistent performance difference between DEET and icaridin in field studies and concluded that they are equally preferred mosquito repellents, noting that 50% DEET offers longer protection but is not available in some countries.
Concentrations
The concentration of DEET in products may range from less than 10% to nearly 100%, but concentrations greater than 50% do not increase the duration of protection. Higher concentrations can be safely applied to clothing, although it may damage some types of synthetic fibers. In the United Kingdom, the publicly-funded healthcare system, the National Health Service (NHS), recommends that UK citizens should use a concentration of 50% when visiting areas of the world with malaria. A lower concentration of 10% is recommended for infants and children. Health Canada decided to limit DEET concentration to 30% in the country since 2002 due to an increased long-term risk observed with repeated applications.
DEET is often sold and used in spray or lotion in concentrations up to 100%. Consumer Reports found a correlation between DEET concentration and hours of protection against insect bites. 100% DEET was found to offer up to 12 hours of protection while several lower concentration DEET formulations (20–34%) offered 3–6 hours of protection. Other research has corroborated the effectiveness of DEET. The Centers for Disease Control and Prevention (CDC) recommends 30–50% DEET to prevent the spread of pathogens carried by insects.
A 2008 study found that higher concentrations of DEET have an improved ability to repel insects through fabric.
Contraindications
DEET should not be used on children younger than 2 months of age.
Adverse effects
When used as directed, products containing between 10% and 30% DEET have been found by the American Academy of Pediatrics to be safe to use on children as well as adults.
As a precaution, manufacturers advise that DEET products should not be used under clothing or on damaged skin, and that preparations be washed off after they are no longer needed or between applications. DEET can irritate the eyes and, unlike icaridin, it can cause breathing difficulty, headaches, or, in rare cases, it may cause severe epidermal reactions.
The authors of a 2002 study published in The New England Journal of Medicine wrote:
... this repellent has been subjected to more scientific and toxicologic scrutiny than any other repellent substance. ... DEET has a remarkable safety profile after 40 years of use and nearly 8 billion human applications. Fewer than 50 cases of serious toxic effects have been documented in ... medical literature since 1960 ... Many of these cases of toxic effects involved long-term, heavy, frequent, or whole-body application of DEET. No correlation has been found between the concentration of DEET used and the risk of toxic effects. ... When applied with common sense, DEET-based repellents can be expected to provide a safe as well as a long-lasting repellent effect ... under circumstances in which it is crucial to be protected against arthropod bites that might transmit disease.
In the DEET Reregistration Eligibility Decision (RED) in 1998, the United States Environmental Protection Agency (EPA) reported 14 to 46 cases of potential DEET-associated seizures, including four deaths. The EPA states: "... it does appear that some cases are likely related to DEET toxicity," which may underreport the risk as physicians may fail to check for history of DEET use or fail to report cases of seizure subsequent to DEET use.
In 1997, the Pesticide Information Project of Cooperative Extension Offices of Cornell University stated that "Everglades National Park employees having extensive DEET exposure were more likely to have insomnia, mood disturbances and impaired cognitive function than lesser exposed co-workers".
Citing human health reasons, Health Canada barred the sale of insect repellents for human use that contained more than 30% DEET in a 2002 re-evaluation "based on a human health risk assessment that considered daily application of DEET over a prolonged period of time". The agency recommended that DEET-based products be used on children between the ages of 2 and 12 only if the concentration of DEET is 10% or less and that repellents be applied no more than 3 times a day, children under 2 should not receive more than 1 application of repellent in a day and DEET-based products of any concentration should not be used on infants under 6 months.
A 2020 study performed by students within the University of Florida's College of Public Health and Health Professions analyzed data from the National Health and Nutrition Examination Survey and identified 1,205 participants who had "DEET metabolic levels recorded at or above detection limits". They analyzed biomarkers related to systemic inflammation, immune, liver, and kidney functions, and found no "evidence that DEET exposure has any impact on the biomarkers identified."
Detection in body fluids
DEET may be measured in blood, plasma, or urine by gas or liquid chromatography-mass spectrometry to confirm a diagnosis of poisoning in hospitalized patients or to provide evidence in a medicolegal death investigation. Blood or plasma DEET concentrations are expected to be in a range of 0.3–3.0 mg/L during the first 8 hours after dermal application in persons using the chemical appropriately, >6 mg/L in intoxicated patients and >100 mg/L in victims of acute intentional oral overdose.
Overdose
Applying DEET to the skin is safe if done as directed. Adverse reactions are very rare, about 1 in 100 million people. However, repeated use of DEET in very high concentrations can lead to toxic encephalopathy with severe neurological symptoms including seizures, tremors and slurred speech. The risk is higher for children since they have a greater surface area to body weight ratio.
Interactions
Limited data indicates that combining insect repellents with DEET and sunscreen decreases the sun protection factor of the sunscreen by about a third. Unlike icaridin, the combination also increases the absorption of both significantly. When the two need to be used together, the repellent should be applied after the sunscreen has been absorbed, about 30 or more minutes later.
When DEET is used in combination with insecticides for cockroaches it can strengthen the toxicity of carbamate, an acetylcholinesterase inhibitor. These 1996 findings indicate that DEET has neurological effects on insects in addition to known olfactory effects, and that its toxicity is strengthened in combination with other insecticides.
Damage to materials
Unlike icaridin, DEET is an effective solvent and may dissolve some watch crystals, plastics, rayon, spandex, other synthetic fabrics, and painted or varnished surfaces including nail polish. It also may act as a plasticizer by remaining inside some formerly hard plastics, leaving them softened and more flexible. DEET is incompatible with rayon, acetate, or dynel clothing.
Environmental impact
Though DEET is not expected to bioaccumulate, it has been found to have a slight toxicity for fresh-water fish such as rainbow trout and tilapia, and it also has been shown to be toxic for some species of freshwater zooplankton. DEET has been detected at low concentrations in water bodies as a result of production and use, such as in the Mississippi River and its tributaries, where a 1991 study detected levels varying from 5 to 201 ng/L.
A 1975 study analyzed the effects of DEET on communities of freshwater organisms native to Chinese waterways and found that DEET was moderately toxic to aquatic organisms compared to other commercial insect repellants. The most-at-risk organisms were algae colonies which often experienced "significant biomass decline and community composition shift[s]" when exposed to DEET at 500 ng/L.
DEET is biodegraded by fungi into products less toxic to zooplankton. It degrades well under aerobic conditions, but poorly and slowly under anaerobic conditions.
Mechanism of action
DEET is thought to provide protection from mosquitos via two pathways, both by negatively impacting mosquito odorant receptors at a distance, and by negatively impacting mosquito chemoreceptors upon contact. The exact mechanisms are still being researched, but the two most likely hypotheses are the "smell and avoid hypothesis" (that DEET has an unpleasant odor to insects), and the "bewilderment hypothesis" (that smelling DEET confuses insects). An alternative hypothesis is that DEET "masks" humans by reducing the volatility of skin odorants that are attractive to insects.
Synthesis
A slightly yellow liquid at room temperature, it can be prepared by converting m-toluic acid (3-methylbenzoic acid) to the corresponding acyl chloride using thionyl chloride (SOCl2), and then allowing that product to react with diethylamine:
History
DEET was developed in 1944 by Samuel Gertler of the United States Department of Agriculture for use by the United States Army, following its experience of jungle warfare during World War II. It was originally tested as a pesticide on farm fields, and entered military use in 1946 and civilian use in 1957. It was used in Vietnam and Southeast Asia.
In its original form, known as "bug juice", the application solution was composed of 75% DEET and 25% ethanol. Later, a new version of the repellent was developed by the U.S. Army and the USDA. This formulation consisted of DEET and a mixture of polymers that extended its release and reduced its evaporation rate. This extended-release application was registered by the Environmental Protection Agency in 1991.
| Physical sciences | Amides and amines | Chemistry |
345938 | https://en.wikipedia.org/wiki/Libellulidae | Libellulidae | The chasers, darters, skimmers, and perchers and their relatives form the Libellulidae, the largest family of dragonflies. It is sometimes considered to contain the Corduliidae as the subfamily Corduliinae and the Macromiidae as the subfamily Macromiinae. Even if these are excluded (as Silsby does), there remains a family of over 1000 species. With nearly worldwide distribution, these are the most commonly encountered dragonflies.
The genus Libellula is mostly New World but also has one of the few endangered odonates from Japan: Libellula angelina. Many of the members of this genus are brightly colored or have banded wings. The related genus Plathemis includes the whitetails. The genus Celithemis contains several brightly marked species in the southern United States. Members of the genus Sympetrum are called darters (or meadowhawks in North America) and are found throughout most of the world, except Australia. Several tropical species in the genera Trithemis and Zenithoptera are considered to be especially beautiful. Other common genera include Tramea and Pantala.
Libellulids have stout-bodied larvae with the lower lip or labium developed into a mask over the lower part of the face.
Etymology
The family name may have been derived from the Latin libella which means "booklet".
Genera
The Libelluidae contain these genera:
Gallery
| Biology and health sciences | Odonata | Animals |
346030 | https://en.wikipedia.org/wiki/Improper%20integral | Improper integral | In mathematical analysis, an improper integral is an extension of the notion of a definite integral to cases that violate the usual assumptions for that kind of integral. In the context of Riemann integrals (or, equivalently, Darboux integrals), this typically involves unboundedness, either of the set over which the integral is taken or of the integrand (the function being integrated), or both. It may also involve bounded but not closed sets or bounded but not continuous functions. While an improper integral is typically written symbolically just like a standard definite integral, it actually represents a limit of a definite integral or a sum of such limits; thus improper integrals are said to converge or diverge. If a regular definite integral (which may retronymically be called a proper integral) is worked out as if it is improper, the same answer will result.
In the simplest case of a real-valued function of a single variable integrated in the sense of Riemann (or Darboux) over a single interval, improper integrals may be in any of the following forms:
, where is undefined or discontinuous somewhere on
The first three forms are improper because the integrals are taken over an unbounded interval. (They may be improper for other reasons, as well, as explained below.) Such an integral is sometimes described as being of the "first" type or kind if the integrand otherwise satisfies the assumptions of integration. Integrals in the fourth form that are improper because has a vertical asymptote somewhere on the interval may be described as being of the "second" type or kind. Integrals that combine aspects of both types are sometimes described as being of the "third" type or kind.
In each case above, the improper integral must be rewritten using one or more limits, depending on what is causing the integral to be improper. For example, in case 1, if is continuous on the entire interval , then
The limit on the right is taken to be the definition of the integral notation on the left.
If is only continuous on and not at itself, then typically this is rewritten as
for any choice of . Here both limits must converge to a finite value for the improper integral to be said to converge. This requirement avoids the ambiguous case of adding positive and negative infinities (i.e., the "" indeterminate form). Alternatively, an iterated limit could be used or a single limit based on the Cauchy principal value.
If is continuous on and , with a discontinuity of any kind at , then
for any choice of . The previous remarks about indeterminate forms, iterated limits, and the Cauchy principal value also apply here.
The function can have more discontinuities, in which case even more limits would be required (or a more complicated principal value expression).
Cases 2–4 are handled similarly. See the examples below.
Improper integrals can also be evaluated in the context of complex numbers, in higher dimensions, and in other theoretical frameworks such as Lebesgue integration or Henstock–Kurzweil integration. Integrals that are considered improper in one framework may not be in others.
Examples
The original definition of the Riemann integral does not apply to a function such as on the interval , because in this case the domain of integration is unbounded. However, the Riemann integral can often be extended by continuity, by defining the improper integral instead as a limit
The narrow definition of the Riemann integral also does not cover the function on the interval . The problem here is that the integrand is unbounded in the domain of integration. In other words, the definition of the Riemann integral requires that both the domain of integration and the integrand be bounded. However, the improper integral does exist if understood as the limit
Sometimes integrals may have two singularities where they are improper. Consider, for example, the function integrated from 0 to (shown right). At the lower bound of the integration domain, as goes to 0 the function goes to , and the upper bound is itself , though the function goes to 0. Thus this is a doubly improper integral. Integrated, say, from 1 to 3, an ordinary Riemann sum suffices to produce a result of /6. To integrate from 1 to , a Riemann sum is not possible. However, any finite upper bound, say (with ), gives a well-defined result, . This has a finite limit as goes to infinity, namely /2. Similarly, the integral from 1/3 to 1 allows a Riemann sum as well, coincidentally again producing /6. Replacing 1/3 by an arbitrary positive value (with ) is equally safe, giving . This, too, has a finite limit as goes to zero, namely /2. Combining the limits of the two fragments, the result of this improper integral is
This process does not guarantee success; a limit might fail to exist, or might be infinite. For example, over the bounded interval from 0 to 1 the integral of does not converge; and over the unbounded interval from 1 to the integral of does not converge.
It might also happen that an integrand is unbounded near an interior point, in which case the integral must be split at that point. For the integral as a whole to converge, the limit integrals on both sides must exist and must be bounded. For example:
But the similar integral
cannot be assigned a value in this way, as the integrals above and below zero in the integral domain do not independently converge. (However, see Cauchy principal value.)
Convergence of the integral
An improper integral converges if the limit defining it exists. Thus for example one says that the improper integral
exists and is equal to L if the integrals under the limit exist for all sufficiently large t, and the value of the limit is equal to L.
It is also possible for an improper integral to diverge to infinity. In that case, one may assign the value of ∞ (or −∞) to the integral. For instance
However, other improper integrals may simply diverge in no particular direction, such as
which does not exist, even as an extended real number. This is called divergence by oscillation.
A limitation of the technique of improper integration is that the limit must be taken with respect to one endpoint at a time. Thus, for instance, an improper integral of the form
can be defined by taking two separate limits; to which
provided the double limit is finite. It can also be defined as a pair of distinct improper integrals of the first kind:
where c is any convenient point at which to start the integration. This definition also applies when one of these integrals is infinite, or both if they have the same sign.
An example of an improper integral where both endpoints are infinite is the Gaussian integral An example which evaluates to infinity is But one cannot even define other integrals of this kind unambiguously, such as since the double limit is infinite and the two-integral method
yields an indeterminate form, In this case, one can however define an improper integral in the sense of Cauchy principal value:
The questions one must address in determining an improper integral are:
Does the limit exist?
Can the limit be computed?
The first question is an issue of mathematical analysis. The second one can be addressed by calculus techniques, but also in some cases by contour integration, Fourier transforms and other more advanced methods.
Types of integrals
There is more than one theory of integration. From the point of view of calculus, the Riemann integral theory is usually assumed as the default theory. In using improper integrals, it can matter which integration theory is in play.
For the Riemann integral (or the Darboux integral, which is equivalent to it), improper integration is necessary both for unbounded intervals (since one cannot divide the interval into finitely many subintervals of finite length) and for unbounded functions with finite integral (since, supposing it is unbounded above, then the upper integral will be infinite, but the lower integral will be finite).
The Lebesgue integral deals differently with unbounded domains and unbounded functions, so that often an integral which only exists as an improper Riemann integral will exist as a (proper) Lebesgue integral, such as . On the other hand, there are also integrals that have an improper Riemann integral but do not have a (proper) Lebesgue integral, such as . The Lebesgue theory does not see this as a deficiency: from the point of view of measure theory, and cannot be defined satisfactorily. In some situations, however, it may be convenient to employ improper Lebesgue integrals as is the case, for instance, when defining the Cauchy principal value. The Lebesgue integral is more or less essential in the theoretical treatment of the Fourier transform, with pervasive use of integrals over the whole real line.
For the Henstock–Kurzweil integral, improper integration is not necessary, and this is seen as a strength of the theory: it encompasses all Lebesgue integrable and improper Riemann integrable functions.
Improper Riemann integrals and Lebesgue integrals
In some cases, the integral
can be defined as an integral (a Lebesgue integral, for instance) without reference to the limit
but cannot otherwise be conveniently computed. This often happens when the function f being integrated from a to c has a vertical asymptote at c, or if c = ∞ (see Figures 1 and 2). In such cases, the improper Riemann integral allows one to calculate the Lebesgue integral of the function. Specifically, the following theorem holds :
If a function f is Riemann integrable on [a,b] for every b ≥ a, and the partial integrals
are bounded as b → ∞, then the improper Riemann integrals
both exist. Furthermore, f is Lebesgue integrable on [a, ∞), and its Lebesgue integral is equal to its improper Riemann integral.
For example, the integral
can be interpreted alternatively as the improper integral
or it may be interpreted instead as a Lebesgue integral over the set (0, ∞). Since both of these kinds of integral agree, one is free to choose the first method to calculate the value of the integral, even if one ultimately wishes to regard it as a Lebesgue integral. Thus improper integrals are clearly useful tools for obtaining the actual values of integrals.
In other cases, however, a Lebesgue integral between finite endpoints may not even be defined, because the integrals of the positive and negative parts of f are both infinite, but the improper Riemann integral may still exist. Such cases are "properly improper" integrals, i.e. their values cannot be defined except as such limits. For example,
cannot be interpreted as a Lebesgue integral, since
But is nevertheless integrable between any two finite endpoints, and its integral between 0 and ∞ is usually understood as the limit of the integral:
Singularities
One can speak of the singularities of an improper integral, meaning those points of the extended real number line at which limits are used.
Cauchy principal value
Consider the difference in values of two limits:
The former is the Cauchy principal value of the otherwise ill-defined expression
Similarly, we have
but
The former is the principal value of the otherwise ill-defined expression
All of the above limits are cases of the indeterminate form .
These pathologies do not affect "Lebesgue-integrable" functions, that is, functions the integrals of whose absolute values are finite.
Summability
An improper integral may diverge in the sense that the limit defining it may not exist. In this case, there are more sophisticated definitions of the limit which can produce a convergent value for the improper integral. These are called summability methods.
One summability method, popular in Fourier analysis, is that of Cesàro summation. The integral
is Cesàro summable (C, α) if
exists and is finite . The value of this limit, should it exist, is the (C, α) sum of the integral.
An integral is (C, 0) summable precisely when it exists as an improper integral. However, there are integrals which are (C, α) summable for α > 0 which fail to converge as improper integrals (in the sense of Riemann or Lebesgue). One example is the integral
which fails to exist as an improper integral, but is (C,α) summable for every α > 0. This is an integral version of Grandi's series.
Multivariable improper integrals
The improper integral can also be defined for functions of several variables. The definition is slightly different, depending on whether one requires integrating over an unbounded domain, such as , or is integrating a function with singularities, like .
Improper integrals over arbitrary domains
If is a non-negative function that is Riemann integrable over every compact cube of the form , for , then the improper integral of f over is defined to be the limit
provided it exists.
A function on an arbitrary domain A in is extended to a function on by zero outside of A:
The Riemann integral of a function over a bounded domain A is then defined as the integral of the extended function over a cube containing A:
More generally, if A is unbounded, then the improper Riemann integral over an arbitrary domain in is defined as the limit:
Improper integrals with singularities
If f is a non-negative function which is unbounded in a domain A, then the improper integral of f is defined by truncating f at some cutoff M, integrating the resulting function, and then taking the limit as M tends to infinity. That is for , set . Then define
provided this limit exists.
Functions with both positive and negative values
These definitions apply for functions that are non-negative. A more general function f can be decomposed as a difference of its positive part and negative part , so
with and both non-negative functions. The function f has an improper Riemann integral if each of and has one, in which case the value of that improper integral is defined by
In order to exist in this sense, the improper integral necessarily converges absolutely, since
| Mathematics | Integral calculus | null |
346133 | https://en.wikipedia.org/wiki/Annihilation | Annihilation | In particle physics, annihilation is the process that occurs when a subatomic particle collides with its respective antiparticle to produce other particles, such as an electron colliding with a positron to produce two photons. The total energy and momentum of the initial pair are conserved in the process and distributed among a set of other particles in the final state. Antiparticles have exactly opposite additive quantum numbers from particles, so the sums of all quantum numbers of such an original pair are zero. Hence, any set of particles may be produced whose total quantum numbers are also zero as long as conservation of energy, conservation of momentum, and conservation of spin are obeyed.
During a low-energy annihilation, photon production is favored, since these particles have no mass. High-energy particle colliders produce annihilations where a wide variety of exotic heavy particles are created.
The word "annihilation" takes its use informally for the interaction of two particles that are not mutual antiparticles not charge conjugate. Some quantum numbers may then not sum to zero in the initial state, but conserve with the same totals in the final state. An example is the "annihilation" of a high-energy electron antineutrino with an electron to produce a W boson.
If the annihilating particles are composite, such as mesons or baryons, then several different particles are typically produced in the final state.
The inverse of annihilation is pair production, the process in which a high-energy photon converts its energy into mass.
Production of a single boson
If the initial two particles are elementary (not composite), then they may combine to produce only a single elementary boson, such as a photon (), gluon (), , or a Higgs boson (). If the total energy in the center-of-momentum frame is equal to the rest mass of a real boson (which is impossible for a massless boson such as the ), then that created particle will continue to exist until it decays according to its lifetime. Otherwise, the process is understood as the initial creation of a boson that is virtual, which immediately converts into a real particle + antiparticle pair. This is called an s-channel process. An example is the annihilation of an electron with a positron to produce a virtual photon, which converts into a muon and anti-muon. If the energy is large enough, a could replace the photon.
Examples
Electron–positron annihilation
+ → +
When a low-energy electron annihilates a low-energy positron (antielectron), the most probable result is the creation of two or more photons, since the only other final-state Standard Model particles that electrons and positrons carry enough mass–energy to produce are neutrinos, which are approximately 10,000 times less likely to produce, and the creation of only one photon is forbidden by momentum conservation—a single photon would carry nonzero momentum in any frame, including the center-of-momentum frame where the total momentum vanishes. Both the annihilating electron and positron particles have a rest energy of about 0.511 million electron-volts (MeV). If their kinetic energies are relatively negligible, this total rest energy appears as the photon energy of the photons produced. Each of the photons then has an energy of about 0.511 MeV. Momentum and energy are both conserved, with 1.022 MeV of photon energy (accounting for the rest energy of the particles) moving in opposite directions (accounting for the total zero momentum of the system).
If one or both charged particles carry a larger amount of kinetic energy, various other particles can be produced. Furthermore, the annihilation (or decay) of an electron–positron pair into a single photon can occur in the presence of a third charged particle, to which the excess momentum can be transferred by a virtual photon from the electron or positron. The inverse process, pair production by a single real photon, is also possible in the electromagnetic field of a third particle.
Proton–antiproton annihilation
When a proton encounters its antiparticle (and more generally, if any species of baryon encounters the corresponding antibaryon), the reaction is not as simple as electron–positron annihilation. Unlike an electron, a proton is a composite particle consisting of three "valence quarks" and an indeterminate number of "sea quarks" bound by gluons. Thus, when a proton encounters an antiproton, one of its quarks, usually a constituent valence quark, may annihilate with an antiquark (which more rarely could be a sea quark) to produce a gluon, after which the gluon together with the remaining quarks, antiquarks, and gluons will undergo a complex process of rearrangement (called hadronization or fragmentation) into a number of mesons, (mostly pions and kaons), which will share the total energy and momentum. The newly created mesons are unstable, and unless they encounter and interact with some other material, they will decay in a series of reactions that ultimately produce only photons, electrons, positrons, and neutrinos. This type of reaction will occur between any baryon (particle consisting of three quarks) and any antibaryon consisting of three antiquarks, one of which corresponds to a quark in the baryon. (This reaction is unlikely if at least one among the baryon and anti-baryon is exotic enough that they share no constituent quark flavors.) Antiprotons can and do annihilate with neutrons, and likewise antineutrons can annihilate with protons, as discussed below.
Reactions in which proton–antiproton annihilation produces as many as 9 mesons have been observed, while production of 13 mesons is theoretically possible. The generated mesons leave the site of the annihilation at moderate fractions of the speed of light and decay with whatever lifetime is appropriate for their type of meson.
Similar reactions will occur when an antinucleon annihilates within a more complex atomic nucleus, save that the resulting mesons, being strongly interacting, have a significant probability of being absorbed by one of the remaining "spectator" nucleons rather than escaping. Since the absorbed energy can be as much as ~2 GeV, it can in principle exceed the binding energy of even the heaviest nuclei. Thus, when an antiproton annihilates inside a heavy nucleus such as uranium or plutonium, partial or complete disruption of the nucleus can occur, releasing large numbers of fast neutrons. Such reactions open the possibility for triggering a significant number of secondary fission reactions in a subcritical mass and may potentially be useful for spacecraft propulsion.
Higgs production
In collisions of two nucleons at very high energies, sea quarks and gluons tend to dominate the interaction rate, so neither nucleon need be an anti-particle for annihilation of a quark pair or "fusion" of two gluons to occur. Examples of such processes contribute to the production of the long-sought Higgs boson. The Higgs is directly produced very weakly by annihilation of light (valence) quarks, but heavy or sea or produced quarks are available. In 2012, the CERN laboratory in Geneva announced the discovery of the Higgs in the debris from proton–proton collisions at the Large Hadron Collider (LHC). The strongest Higgs yield is from fusion of two gluons (via annihilation of a heavy quark pair), while two quarks or antiquarks produce more easily identified events through radiation of a Higgs by a produced virtual vector boson or annihilation of two such vector bosons.
| Physical sciences | Particle physics: General | Physics |
346312 | https://en.wikipedia.org/wiki/Tachinidae | Tachinidae | The Tachinidae are a large and variable family of true flies within the insect order Diptera, with more than 8,200 known species and many more to be discovered. Over 1,300 species have been described in North America alone. Insects in this family commonly are called tachinid flies or simply tachinids. As far as is known, they all are protelean parasitoids, or occasionally parasites, of arthropods, usually other insects. The family is known from many habitats in all zoogeographical regions and is especially diverse in South America.
Life cycle
Reproductive strategies vary greatly between tachinid species, largely, but not always clearly, according to their respective life cycles. Many species are generalists rather than specialists. Comparatively few are restricted to a single host species, so there is little tendency towards the close co-evolution one finds in the adaptations of many specialist species to their hosts, such as are typical of protelean parasitoids among the Hymenoptera.
Larvae (maggots) of most members of this family are parasitoids (developing inside a living host, ultimately killing it). In contrast, a few are parasitic (not generally killing the host). Tachinid larvae feed on the host tissues, either after having been injected into the host by the parent, or penetrating the host from outside. Various species have different modes of oviposition and of host invasion. Typically, tachinid larvae are endoparasites (internal parasites) of caterpillars of butterflies and moths, or the eruciform larvae of sawflies. For example, they have been found to lay eggs in African sugarcane borer larva, a species of moth common in sub-Saharan Africa, as well as the more northerly Arctic woolly bear moth. However, some species attack adult beetles and some attack beetle larvae. Others attack various types of true bugs, and others attack grasshoppers; a few even attack centipedes. Also parasitised are bees, wasps and sawflies.
Oviposition
Probably the majority of female tachinids lay white, ovoid eggs with flat undersides onto the skin of the host insect. Imms mentions the genera Gymnosoma, Thrixion, Winthemia, and Eutachina as examples. In a closely related strategy some genera are effectively ovoviviparous (some authorities prefer the term ovolarviparous) and deposit a hatching larva onto the host. For example, this occurs in Tachinidae species which parasitize the butterfly Danaus chrysippus in Ghana. The free larvae immediately bore into the host's body. Illustrative genera include Exorista and Voria. Many tachinid eggs hatch quickly, having partly developed inside the mother's uterus, which is long and often coiled for retaining developing eggs. However, it is suggested that the primitive state probably is to stick unembryonated eggs to the surface of the host.
Many other species inject eggs into the host's body, using the extensible, penetrating part of their ovipositor, sometimes called the oviscapt, which roughly translates to "egg digger". Species in the genera Ocyptera, Alophora, and Compsilura are examples.
In many species only one egg is laid on or in any individual host, and accordingly such an egg tends to be large, as is typical for eggs laid in small numbers. They are large enough to be clearly visible if stuck onto the outside of the host, and they generally are so firmly stuck that eggs cannot be removed from the skin of the host without killing them. Furthermore, scientists have observed in studies with the host cabbage looper that being glued to the host insect helps maggots burrow into the larva, where they remain until fully developed.
Yet another strategy of oviposition among some Tachinidae is to lay large numbers of small, darkly coloured eggs on the food plants of the host species. Sturmia, Zenillia, and Gonia are such genera.
Many tachinids are important natural enemies of major insect pests, and some species actually are used in biological pest control; for example, some species of tachinid flies have been introduced into North America from their native lands as biocontrols to suppress populations of alien pests. Conversely, certain tachinid flies that prey on useful insects are themselves considered as pests; they can present troublesome problems in the sericulture industry by attacking silkworm larvae. One particularly notorious silkworm pest is the Uzi fly (Exorista bombycis).
Another reproductive strategy is to leave the eggs in the host's environment; for example, the female might lay on leaves, where the host is likely to ingest them. Some tachinids that are parasitoids of stem-boring caterpillars deposit eggs outside the host's burrow, letting the first instar larvae do the work of finding the host for themselves. In other species, the maggots use an ambush technique, waiting for the host to pass and then attacking it and burrowing into its body.
Adult tachinids are not parasitic, but either do not feed at all or visit flowers, decaying matter, or similar sources of energy to sustain themselves until they have concluded their procreative activities. Their non-parasitic behaviour after eclosion from the pupa is what justifies the application of the term "protelean".
Description
Tachinid flies are extremely varied in appearance. Some adult flies may be brilliantly colored and resemble blow-flies (family Calliphoridae). Most however are rather drab, some resembling house flies. However, tachinid flies commonly are more bristly and more robust. Also, they usually have a characteristic appearance. They have three-segmented antennae, a diagnostically prominent postscutellum bulging beneath the scutellum (a segment of the mesonotum). They are aristate flies, and the arista usually is bare, though sometimes plumose. The calypters (small flaps above the halteres) are usually very large. Their fourth long vein bends away sharply.
Adult flies feed on flowers and nectar from aphids and scale insects. As many species typically feed on pollen, they can be important pollinators of some plants, especially at higher elevations in mountains where bees are relatively few.
The taxonomy of this family presents many difficulties. It is largely based on morphological characters of the adult flies, but also on reproductive habits and on the immature stage.
As biological pest control
Some tachinid flies parasitize pest species. This has allowed them to be used as biological control agents by farmers. Some Tachinidae are generalists; for instance, Compsilura concinnata uses, at least, 200 different hosts, and thus are less safe to be used as biological controls because they will attack non-pest species, resulting in population decline. Others are more specialized and are safer; for instance, Istocheta aldrichi, which only attacks the Japanese beetle.
Evolution
This clade appears to have originated in the middle Eocene. The oldest known putatively tachinid fossil (Lithexorista) dates from the Eocene Green River Formation in Wyoming.
| Biology and health sciences | Flies (Diptera) | Animals |
346315 | https://en.wikipedia.org/wiki/Human%20variability | Human variability | Human variability, or human variation, is the range of possible values for any characteristic, physical or mental, of human beings.
Frequently debated areas of variability include cognitive ability, personality, physical appearance (body shape, skin color, etc.) and immunology.
Variability is partly heritable and partly acquired (nature vs. nurture debate).
As the human species exhibits sexual dimorphism, many traits show significant variation not just between populations but also between the sexes.
Sources of human variability
Human variability is attributed to a combination of environmental and genetic sources including:
For the genetic variables listed above, few of the traits characterizing human variability are controlled by simple Mendelian inheritance. Most are polygenic or are determined by a complex combination of genetics and environment.
Many genetic differences (polymorphisms) have little effect on health or reproductive success but help to distinguish one population from another. It is helpful for researchers in the field of population genetics to study ancient migrations and relationships between population groups.
Environmental factors
Climate and disease
Other important factors of environmental factors include climate and disease. Climate has effects on determining what kinds of human variation are more adaptable to survive without much restrictions and hardships. For example, people who live in a climate where there is a lot of exposure to sunlight have a darker color of skin tone. Evolution has caused production of folate (folic acid) from UV radiation, thus giving them darker skin tone with more melanin to make sure child development is smooth and successful. Conversely, people who live farther away from the equator have a lighter skin tone. This is due to a need for an increased exposure and absorbance of sunlight to make sure the body can produce enough vitamin D for survival.
Blackfoot disease is a disease caused by environmental pollution and causes people to have black, charcoal-like skin in the lower limbs. This is caused by arsenic pollution in water and food source. This is an example of how disease can affect human variation. Another disease that can affect human variation is syphilis, a sexual transmitted disease. Syphilis does not affect human variation until the middle stage of the disease. It then starts to grow rashes all over the body, affecting people's human variation.
Nutrition
Phenotypic variation is a combination of one's genetics and their surrounding environment, with no interaction or mutual influence between the two. This means that a significant portion of human variability can be controlled by human behavior. Nutrition and diet play a substantial role in determining phenotype because they are arguably the most controllable forms of environmental factors that create epigenetic changes. This is because they can be changed or altered relatively easily as opposed to other environmental factors like location.
If people are reluctant to changing their diets, consuming harmful foods can have chronic negative effects on variability. One such instance of this occurs when eating certain chemicals through one's diet or consuming carcinogens, which can have adverse effects on individual phenotype. For example, Bisphenol A (BPA) is a known endocrine disruptor that mimics the hormone estradiol and can be found in various plastic products. BPA seeps into food or drinks when the plastic containing it is heated up and begins to melt. When these contaminated substances are consumed, especially often and over long periods of time, one's risk of diabetes and cardiovascular disease increases. BPA also has the potential to alter "physiological weight control patterns." Examples such as this demonstrate that preserving a healthy phenotype largely rests on nutritional decision-making skills.
The concept that nutrition and diet affect phenotype extends to what the mother eats during pregnancy, which can have drastic effects on the outcome of the phenotype of the child. A recent study by researchers at the MRC International Nutrition Group shows that "methylation machinery can be disrupted by nutrient deficiencies and that this can lead to disease" susceptibility in newborn babies. The reason for this is because methyl groups have the ability to silence certain genes. Increased deficiencies of various nutrients such as this have the potential to permanently change the epigenetics of the baby.
Genetic factors
Genetic variation in humans may mean any variance in phenotype which results from heritable allele expression, mutations, and epigenetic changes. While human phenotypes may seem diverse, individuals actually differ by only 1 in every 1,000 base pairs and is primarily the result of inherited genetic differences. Pure consideration of alleles is often referred to as Mendelian Genetics, or more properly Classical Genetics, and involves the assessment of whether a given trait is dominant or recessive and thus, at what rates it will be inherited. The color of one's eyes was long believed to occur with a pattern of brown-eye dominance, with blue eyes being a recessive characteristic resulting from a past mutation. However, it is now understood that eye color is controlled by various genes, and thus, may not follow as distinct a pattern as previously believed. The trait is still the result of variance in genetic sequence between individuals as a result of inheritance from their parents. Common traits which may be linked to genetic patterns are earlobe attachment, hair color, and hair growth patterns.
In terms of evolution, genetic mutations are the origins of differences in alleles between individuals. However, mutations may also occur within a person's life-time and be passed down from parent to offspring. In some cases, mutations may result in genetic diseases, such as Cystic Fibrosis, which is the result of a mutation to the CFTR gene that is recessively inherited from both parents. In other cases, mutations may be harmless or phenotypically unnoticeable. We are able to treat biological traits as manifestations of either a single loci or multiple loci, labeling said biological traits as either monogenic or polygenic, respectively. Concerning polygenic traits it may be essential to be mindful of inter-genetic interactions or epistasis. Although epistasis is a significant genetic source of biological variation, it is only additive interactions that are heritable as other epistatic interactions involve recondite inter-genetic relationships. Epistatic interactions in of themselves vary further with their dependency on the results of the mechanisms of recombination and crossing over.
The ability of genes to be expressed may also be a source of variation between individuals and result in changes to phenotype. This may be the result of epigenetics, which are founded upon an organism's phenotypic plasticity, with such a plasticity even being heritable. Epigenetics may result from methylation of gene sequences leading to the blocking of expression or changes to histone protein structuring as a result of environmental or biological cues. Such alterations influence how genetic material is handled by the cell and to what extent certain DNA sections are expressed and compose the epigenome. The division between what can be considered as a genetic source of biological variation and not becomes immensely arbitrary as we approach aspects of biological variation such as epigenetics. Indeed, gene specific gene expression and inheritance may be reliant on environmental influences.
Cultural factors
Archaeological findings such as those that indicate that the Middle Stone Age and the Acheulean – identified as a specific 'cultural phases' of humanity with a number of characteristics – lasted substantially longer in some places or 'ended' at times over 100,000 years apart, highlight a significant spatiotemporal cultural variability in and complexity of the sociocultural history and evolution of humanity. In some cases cultural factors may be intertwined with genetic and environmental factors.
Measuring variation
Scientific
Measurement of human variation can fall under the purview of several scholarly disciplines, many of which lie at the intersection of biology and statistics. The methods of biostatistics, the application of statistical methods to the analysis of biological data, and bioinformatics, the application of information technologies to the analysis of biological data, are utilized by researchers in these fields to uncover significant patterns of variability. Some fields of scientific research include the following:
Demography is a branch of statistics and sociology concerned with the statistical study of populations, especially humans. A demographic analysis can measure various metrics of a population, most commonly metrics of size and growth, diversity in culture, ethnicity, language, religious belief, political belief, etc. Biodemography is a subfield which specifically integrates biological understanding into demographics analysis.
In the social sciences, social research is conducted and collected data is analyzed under statistical methods. The methodologies of this research can be divided into qualitative and quantitative designs. Some example subdisciplines include:
Anthropology, the study of human societies. Comparative research in subfields of anthropology may yield results on human variation with respect to the subfield's topic of interest.
Psychology, the study of behavior from a mental perspective. Does a lot of experiments and analysis grouped into quantitative or qualitative research methods.
Sociology, the study of behavior from a social perspective. Sociological research can be conducted in either quantitative or qualitative formats, depending on the nature of data collected and the subfield of sociology under which the research falls. Analysis of this data is subject to quantitative or qualitative methods. Computational sociology is also a method of producing useful data for studies of social behavior.
Anthropometry
Anthropometry is the study of the measurements of different parts of the human body. Common measurements include height, weight, organ size (brain, stomach, penis, vagina), and other bodily metrics such as waist–hip ratio. Each measurement can vary significantly between populations; for instance, the average height of males of European descent is 178 cm ± 7 cm and of females of European descent is 165 cm ± 7 cm. Meanwhile, average height of Nilotic males in Dinka is 181.3 cm.
Applications of anthropometry include ergonomics, biometrics, and forensics. Knowing the distribution of body measurements enable designers to build better tools for workers. Anthropometry is also used when designing safety equipment such as seat belts. In biometrics, measurements of fingerprints and iris patterns can be used for secure identification purposes.
Measuring genetic variation
Human genomics and population genetics are the study of the human genome and variome, respectively. Studies in these areas may concern the patterns and trends in human DNA. The Human Genome Project and The Human Variome Project are examples of large scale studies of the entire human population to collect data which can be analyzed to understand genomic and genetic variation in individuals, respectively.
The Human Genome Project is the largest scientific project in the history of biology. At a cost of $3.8 billion in funding and over a period of 13 years from 1990 to 2003, the project processed through DNA sequencing the approximately 3 billion base pairs and catalogued the 20,000 to 25,000 genes in human DNA. The project made the data available to all scientific researchers and developed analytical tools for processing this information. A particular finding regarding human variability due to difference in DNA made possible by the Human Genome Project is that any two individuals share 99.9% of their nucleotide sequences.
The Human Variome Project is a similar undertaking with the goal of identification and categorization of the set of human genetic variation, specifically variations which are medically pertinent. This project will also provide a data repository for further research and analysis of disease. The Human Variome Project was launched in 2006 and is being run by an international community of researchers and representatives, including collaborators from the World Health Organization and the United Nations Educational, Scientific, and Cultural Organization.
Genetic drift
Genetic drift is one method by which variability occurs in populations. Unlike natural selection, genetic drift occurs when alleles decrease randomly over time and not as a result of selection bias. Over a long history, this can cause significant shifts in the underlying genetic distribution of a population. We can model genetic drift with the Wright-Fisher model. In a population of N with 2N genes, there are two alleles with frequencies p and q. If the previous generation had an allele with frequency p, then the probability that the next generation has k of that allele is:
Over time, one allele will be fixed when the frequency of that allele reaches 1 and the frequency of the other allele reaches 0. The probability that any allele is fixed is proportional to the frequency of that allele. For two alleles with frequencies p and q, the probability that p will be fixed is p. The expected number of generations for an allele with frequency p to be fixed is:
Where Ne is the effective population size.
Single-nucleotide polymorphism
Single-nucleotide polymorphism or SNPs are variations of a single nucleotide. SNPs can occur in coding or non-coding regions of genes and on average occur once every 300 nucleotides. SNPs in coding regions can cause synonymous, missense, and nonsense mutations. SNPs have shown to be correlated with drug responses and risk of diseases such as sickle-cell anemia, Alzheimer's disease, cystic fibrosis, and more.
DNA fingerprinting
DNA profiling, whereby a DNA fingerprint is constructed by extracting a DNA sample from body tissue or fluid. Then, it is segmented using restriction enzymes and each segment marked with probes then exposed on X-ray film. The segments form patterns of black bars;the DNA fingerprint. DNA Fingerprints are used in conjunction with other methods in order to individuals information in Federal programs such as CODIS (Combined DNA Index System for Missing Persons) in order to help identify individuals
Mitochondrial DNA
Mitochondrial DNA, which is only passed from mother to child. The first human population studies based on mitochondrial DNA were performed by restriction enzyme analyses (RFLPs) and revealed differences between the four ethnic groups (Caucasian, Amerindian, African, and Asian). Differences in mtDNA patterns have also been shown in communities with a different geographic origin within the same ethnic group
Alloenzymic variation
Alloenzymic variation, a source of variation that identifies protein variants of the same gene due to amino acid substitutions in proteins. After grinding tissue to release the cytoplasm, wicks are used to absorb the resulting extract and placed in a slit cut into a starch gel. A low current is run across the gel resulting in a positive and negative ends. Proteins are then separated by charge and size, with the smaller and more highly charged molecules moving more quickly across the gel. This techniques does underestimate true genetic variability as there may be an amino acid substitution but if the amino acid is not charged differently than the original no difference in migration will appear it is estimated that approximately 1/3 of the true genetic variation is not expressed by this technique.
Structural variation
Structural variation, which can include insertions, deletions, duplications, and mutations in DNA. Within the human population, about 13% of the human genome is defined as structurally variant.
Phenotypic variation
Phenotypic variation, which accounts for both genetic and epigenetic factors that affect what characteristics are shown. For applications such as organ donations and matching, phenotypic variation of blood type, tissue type, and organ size are considered.
Civic
Measurement of human variation may also be initiated by governmental parties. A government may conduct a census, the systematic recording of an entire population of a region. The data may be used for calculating metrics of demography such as sex, gender, age, education, employment, etc.; this information is utilized for civic, political, economic, industrial, and environmental assessment and planning.
Commercial
Commercial motivation for understanding variation in human populations arises from the competitive advantage of tailoring products and services for a specific target market. A business may undertake some form of market research in order to collect data on customer preference and behavior and implement changes which align with the results.
Social significance and valuation
Both individuals and entire societies and cultures place values on different aspects of human variability; however, values can change as societies and cultures change. Not all people agree on the values or relative rankings, and neither do all societies and cultures. Nonetheless, nearly all human differences have a social value dimension. Examples of variations which may be given different values in different societies include skin color and/or body structure. Race and sex have a strong value difference, while handedness has a much weaker value difference. The values given to different traits among human variability are often influenced by what phenotypes are more prevalent locally. Local valuation may affect social standing, reproductive opportunities, or even survival.
Differences may vary or be distributed in various ways. Some, like height for a given sex, vary in close to a "normal" or Gaussian distribution. Other characteristics (e.g., skin color) vary continuously in a population, but the continuum may be socially divided into a small number of distinct categories. Then, there are some characteristics that vary bimodally (for example, handedness), with fewer people in intermediate categories.
Classification and evaluation of traits
When an inherited difference of body structure or function is severe enough to cause a significant hindrance in certain perceived abilities, it is termed a genetic disease, but even this categorization has fuzzy edges. There are many instances in which the degree of negative value of a human difference depends completely on the social or physical environment. For example, in a society with a large proportion of deaf people (as Martha's Vineyard in the 19th century), it was possible to deny that deafness is a disability. Another example of social renegotiation of the value assigned to a difference is reflected in the controversy over management of ambiguous genitalia, especially whether abnormal genital structure has enough negative consequences to warrant surgical correction.
Furthermore, many genetic traits may be advantageous in certain circumstances and disadvantageous in others. Being a heterozygote or carrier of the sickle-cell disease gene confers some protection against malaria, apparently enough to maintain the gene in populations of malarial areas. In a homozygous dose it is a significant disability.
Each trait has its own advantages and disadvantages, but sometimes a trait that is found desirable may not be favorable in terms of certain biological factors such as reproductive fitness, and traits that are not highly valued by the majority of people may be favorable in terms of biological factors. For example, women tend to have fewer pregnancies on average than before and therefore net worldwide fertility rates are dropping. Moreover, this leads to the fact that multiple births tend to be favorable in terms of number of children and therefore offspring count; when the average number of pregnancies and the average number of children was higher, multiple births made only a slight relative difference in number of children. However, with fewer pregnancies, multiple births can make the difference in number of children relatively large. A hypothetical scenario would be that couple 1 has ten children and couple 2 has eight children, but in both couples, the woman undergoes eight pregnancies. This is not a large difference in ratio of fertility. However, another hypothetical scenario can be that couple 1 has three children and couple 2 has one child but in both couples the woman undergoes one pregnancy (in this case couple 2 has triplets). When the proportion of offspring count in the latter hypothetical scenario is compared, the difference in proportion of offspring count becomes higher. A trait in women known to greatly increase the chance of multiple births is being a tall woman (presumably the chance is further increased when the woman is very tall among both women and men). Yet very tall women are not viewed as a desirable phenotype by the majority of people, and the phenotype of very tall women has not been highly favored in the past. Nevertheless, values placed on traits can change over time.
Such an example is homosexuality. In Ancient Greece, what in present terms would be called homosexuality, primarily between a man and a young boy, was not uncommon and was not outlawed. However, homosexuality became more condemned. Attitudes towards homosexuality alleviated in modern times.
Acknowledgement and study of human differences does have a wide range of uses, such as tailoring the size and shape of manufactured items. See Ergonomics.
Controversies of sociocultural and personal implications
Possession of above average amounts of some abilities is valued by most societies. Some of the traits that societies try to measure by perception are intellectual aptitude in the form of ability to learn, artistic prowess, strength, endurance, agility, and resilience.
Each individual's distinctive differences, even the negatively valued or stigmatized ones, are usually considered an essential part of self-identity.
Membership or status in a social group may depend on having specific values for certain attributes. It is not unusual for people to deliberately try to amplify or exaggerate differences, or to conceal or minimize them, for a variety of reasons. Examples of practices designed to minimize differences include tanning, hair straightening, skin bleaching, plastic surgery, orthodontia, and growth hormone treatment for extreme shortness. Conversely, male-female differences are enhanced and exaggerated in most societies.
In some societies, such as the United States, circumcision is practiced on a majority of males, as well as sex reassignment on intersex infants, with substantial emphasis on cultural and religious norms. Circumcision is highly controversial because although it offers health benefits, such as less chance of urinary tract infections, STDs, and penile cancer, it is considered a drastic procedure that is not medically mandatory and argued as a decision that should be taken when the child is old enough to decide for himself. Similarly, sex reassignment surgery offers psychiatric health benefits to transgender people but is seen as unethical by some Christians, especially when performed on children.
Much controversy surrounds the assigning or distinguishing of some variations, especially since differences between groups in a society or between societies is often debated as part of either a person's "essential" nature or a socially constructed attribution. For example, there has long been a debate among sex researchers on whether sexual orientation is due to evolution and biology (the "essentialist" position), or a result of mutually reinforcing social perceptions and behavioral choices (the "constructivist" perspective). The essentialist position emphasizes inclusive fitness as the reason homosexuality has not been eradicated by natural selection. Gay or lesbian individuals have not been greatly affected by evolutionary selection because they may help the fitness of their siblings and siblings' children, thus increasing their own fitness through inclusive fitness and maintaining evolution of homosexuality. Biological theories for same gender sexual orientation include genetic influences, neuroanatomical factors, and hormone differences but research so far has not provided any conclusive results. In contrast, the social constructivist position argues that sexuality is a result of culture and has originated from language or dialogue about sex. Mating choices are the product of cultural values, such as youth and attractiveness, and homosexuality varies greatly between cultures and societies. In this view, complexities, such as sexual orientation changing during the course of one's lifespan, are accounted for.
Controversy also surrounds the boundaries of "wellness", "wholeness," or "normality." In some cultures, differences in physical appearance, mental ability, and even sex can exclude one from traditions, ceremonies, or other important events, such as religious service. For example, in India, menstruation is not only a taboo subject but also traditionally considered shameful. Depending on beliefs, a woman who is menstruating is not allowed to cook or enter spiritual areas because she is "impure" and "cursed". There has been large-scale renegotiation of the social significance of variations which reduce the ability of a person to do one or more functions in western culture. Laws have been passed to alleviate the reduction of social opportunity available to those with disabilities. The concept of "differently abled" has been pushed by those persuading society to see limited incapacities as a human difference of less negative value.
Ideologies of superiority and inferiority
The extreme exercise of social valuation of human difference is in the definition of "human." Differences between humans can lead to an individual's "nonhuman" status, in the sense of withholding identification, charity, and social participation. Views of these variations can change enormously between cultures over time. For example, nineteenth-century European and American ideas of race and eugenics culminated in the attempts of the Nazi-led German society of the 1930s to deny not just reproduction, but life itself to a variety of people with "differences" attributed in part to biological characteristics. Hitler and Nazi leaders wanted to create a "master race" consisting of only Aryans, or blue-eyed, blonde-haired, and tall individuals, thus discriminating and attempting to exterminate those who didn't fit into this ideal.
Contemporary controversy continues over "what kind of human" is a fetus or child with a significant disability. On one end are people who would argue that Down syndrome is not a disability but a mere "difference," and on the other those who consider it such a calamity as to assume that such a child is better off "not born". For example, in India and China, being female is widely considered such a negatively valued human difference that female infanticide occurs such to severely affect the proportion of sexes.
Common human variations
| Biology and health sciences | Human anatomy | Health |
346451 | https://en.wikipedia.org/wiki/Abrasive | Abrasive | An abrasive is a material, often a mineral, that is used to shape or finish a workpiece through rubbing which leads to part of the workpiece being worn away by friction. While finishing a material often means polishing it to gain a smooth, reflective surface, the process can also involve roughening as in satin, matte or beaded finishes. In short, the ceramics which are used to cut, grind and polish other softer materials are known as abrasives.
Abrasives are extremely commonplace and are used very extensively in a wide variety of industrial, domestic, and technological applications. This gives rise to a large variation in the physical and chemical composition of abrasives as well as the shape of the abrasive. Some common uses for abrasives include grinding, polishing, buffing, honing, cutting, drilling, sharpening, lapping, and sanding (see abrasive machining). (For simplicity, "mineral" in this article will be used loosely to refer to both minerals and mineral-like substances whether man-made or not.)
Files are not abrasives; they remove material not by scratching or rubbing, but by the cutting action of sharp teeth which have been cut into the surface of the file, very much like those of a saw. However, diamond files are a form of coated abrasive (as they are metal rods coated with diamond powder).
Mechanics of abrasion
Abrasives generally rely upon a difference in hardness between the abrasive and the material being worked upon, the abrasive being the harder of the two substances. However, it is not strictly necessary, as any two solid materials that repeatedly rub against each other will tend to wear each other away; examples include, softer shoe soles wearing away wooden or stone steps over decades or centuries or glaciers abrading stone valleys.
Typically, materials used as abrasives are either hard minerals (rated at 7 or above on Mohs scale of mineral hardness) or are synthetic stones, some of which may be chemically and physically identical to naturally occurring minerals but which cannot be called minerals as they did not arise naturally. (While useful for comparative purposes, the Mohs scale is of limited value to materials engineers as it is an arbitrary, ordinal, irregular scale.) Diamond, a common abrasive, for instance occurs both naturally and is industrially produced, as is corundum which occurs naturally but which is nowadays more commonly manufactured from bauxite. However, even softer minerals like calcium carbonate are used as abrasives, such as "polishing agents" in toothpaste.
These minerals are either crushed or are already of a sufficiently small size (anywhere from macroscopic grains as large as about 2 mm to microscopic grains about 0.001 mm in diameter) to permit their use as an abrasive. These grains, commonly called grit, have rough edges, often terminating in points which will decrease the surface area in contact and increase the localised contact pressure. The abrasive and the material to be worked are brought into contact while in relative motion to each other. Force applied through the grains causes fragments of the worked material to break away, while simultaneously smoothing the abrasive grain and/or causing the grain to work loose from the rest of the abrasive.
Some factors which will affect how quickly a substance is abraded include:
Difference in hardness between the two substances: a much harder abrasive will cut faster and deeper
Grain size (grit size): larger grains will cut faster as they also cut deeper
Adhesion between grains, between grains and backing, between grains and matrix: determines how quickly grains are lost from the abrasive and how soon fresh grains, if present, are exposed
Contact force: more force will cause faster abrasion
Loading: worn abrasive and cast off work material tends to fill spaces between abrasive grains so reducing cutting efficiency while increasing friction
Use of lubricant/coolant/metalworking fluid: Can carry away swarf (preventing loading), transport heat (which may affect the physical properties of the workpiece or the abrasive), decrease friction (with the substrate or matrix), suspend worn work material and abrasives allowing for a finer finish, conduct stress to the workpiece.
Abrasive minerals
Abrasives may be classified as either natural or synthetic. When discussing sharpening stones, natural stones have long been considered superior but advances in material technology are seeing this distinction become less distinct. Many synthetic abrasives are effectively identical to a natural mineral, differing only in that the synthetic mineral has been manufactured rather than mined. Impurities in the natural mineral may make it less effective.
Some naturally occurring abrasives are:
Calcite (calcium carbonate)
Emery (impure corundum)
Diamond dust (synthetic diamonds are used extensively)
Novaculite
Pumice
Iron(III) oxide
Sand
Corundum
Garnet
Sandstone
Rotten stone (Tripoli)
Powdered feldspar
Staurolite
Some abrasive minerals (such as zirconia alumina) occur naturally but are sufficiently rare or sufficiently more difficult or costly to obtain such that a synthetic stone is used industrially. These and other artificial abrasives include:
Borazon (cubic boron nitride or CBN)
Ceramic
Ceramic aluminium oxide
Ceramic iron oxide
Corundum (alumina or aluminium oxide)
Dry ice
Glass powder
Steel abrasive
Silicon carbide (carborundum)
Zirconia alumina
Boron carbide
Slags
Manufactured abrasives
Abrasives are shaped for various purposes. Natural abrasives are often sold as dressed stones, usually in the form of a rectangular block. Both natural and synthetic abrasives are commonly available in a wide variety of shapes, often coming as bonded or coated abrasives, including blocks, belts, discs, wheels, sheets, rods and loose grains.
Bonded abrasives
A bonded abrasive is composed of an abrasive material contained within a matrix, although very fine aluminium oxide abrasive may comprise sintered material. This matrix is called a binder and is often a clay, a resin, a glass or a rubber. This mixture of binder and abrasive is typically shaped into blocks, sticks, or wheels. The most common abrasive used is aluminium oxide. Also common are silicon carbide, tungsten carbide and garnet. Artificial sharpening stones are often a bonded abrasive and are readily available as a two sided block, each side being a different grade of grit.
Grinding wheels are cylinders that are rotated at high speed. While once worked with a foot pedal or hand crank, the introduction of electric motors has made it necessary to construct the wheel to withstand greater radial stress to prevent the wheel flying apart as it spins. Similar issues arise with cutting wheels, which are often structurally reinforced with impregnated fibres. High relative speed between abrasive and workpiece often makes necessary the use of a lubricant of some kind. Traditionally, they were called coolants as they were used to prevent frictional heat build up which could damage the workpiece (such as ruining the temper of a blade). Some research suggests that the heat transport property of a lubricant is less important when dealing with metals as the metal will quickly conduct heat from the work surface. More important are their effects upon lessening tensile stresses while increasing some compressive stresses and reducing "thermal and mechanical stresses during chip formation".
Various shapes are also used as heads on rotary tools used in precision work, such as scale modelling.
Bonded abrasives need to be trued and dressed after they are used. Dressing is the cleaning of the waste material (swarf and loose abrasive) from the surface and exposing fresh grit. Depending upon the abrasive and how it was used, dressing may involve the abrasive being simply placed under running water and brushed with a stiff brush for a soft stone or the abrasive being ground against another abrasive, such as aluminium oxide used to dress a grinding wheel.
Truing is restoring the abrasive to its original surface shape. Wheels and stones tend to wear unevenly, leaving the cutting surface no longer flat (said to be "dished out" if it is meant to be a flat stone) or no longer the same diameter across the cutting face. This will lead to uneven abrasion and other difficulties.
Coated abrasives
A coated abrasive comprises an abrasive fixed to a backing material such as paper, cloth, rubber, resin, polyester or even metal, many of which are flexible. Sandpaper is a very common coated abrasive. Coated abrasives are most commonly the same minerals as are used for bonded abrasives. A bonding agent (often some sort of adhesive or resin) is applied to the backing to provide a flat surface to which the grit is then subsequently adhered. A woven backing may also use a filler agent (again, often a resin) to provide additional resilience.
Coated abrasives may be shaped for use in rotary and orbital sanders, for wrapping around sanding blocks, as handpads, as closed loops for use on belt grinders, as striking surfaces on matchboxes, on diamond plates and diamond steels. Diamond tools, though for cutting, are often abrasive in nature.
Other abrasives and their uses
Sand, glass beads, metal pellets copper slag and dry ice may all be used for a process called sandblasting (or similar, such as the use of glass beads which is "bead blasting"). Dry ice will sublimate leaving behind no residual abrasive.
Cutting compound used on automotive paint is an example of an abrasive suspended in a liquid, paste or wax, as are some polishing liquids for silverware and optical media. The liquid, paste or wax acts as a binding agent that keeps the abrasive attached to the cloth which is used as a backing to move the abrasive across the work piece. On cars in particular, wax may serve as both a protective agent by preventing exposure of the paint of metal to air and also act as an optical filler to make scratches less noticeable. Toothpaste contains calcium carbonate or silica as a "polishing agent" to remove plaque and other matter from teeth as the hardness of calcium carbonate is less than that of tooth enamel but more than that of the contaminating agent.
Very fine rouge powder was commonly used for grinding glass, being somewhat replaced by modern ceramics, and is still used in jewellery making for a highly reflective finish.
Cleaning products may also contain abrasives suspended in a paste or cream. They are chosen to be reasonably safe on some linoleum, tile, metal or stone surfaces. However, many laminate surfaces and ceramic topped stoves are easily damaged by these abrasive compounds. Even ceramic/pottery tableware or cookware can damage these surfaces, particularly the bottom of the tableware, which is often unglazed in part or in whole and acts as simply another bonded abrasive.
Metal pots and stoves are often scoured with abrasive cleaners, typically in the form of the aforementioned cream or paste or of steel wool and non woven scouring pads which holds fine grits abrasives.
Human skin is also subjected to abrasion in the form of exfoliation. Abrasives for this can be much softer and more exotic than for other purposes and may include things like almond and oatmeal. Dermabrasion and microdermabrasion are now rather commonplace cosmetic procedures which use mineral abrasives.
Scratched compact discs and DVDs may sometimes be repaired through buffing with a very fine compound, the principle being that a multitude of small scratches will be more optically transparent than a single large scratch. However, this does take some skill and will eventually cause the protective coating of the disc to be entirely eroded (especially if the original scratch is deep), at which time, the data surface will be destroyed if abrasion continues.
Silicon carbide powders are commonly used as abrasive materials in various machining processes, including grinding, water-jet cutting, and sandblasting. These powders are effective for fine grinding or rough polishing of semiconductors, ceramics, and ferrous materials.
Choice of abrasive
The shape, size and nature of the workpiece and the desired finish will influence the choice of the abrasive used. A bonded abrasive grind wheel may be used to commercially sharpen a knife (producing a hollow grind), but an individual may then sharpen the same knife with a natural sharpening stone or an even flexible coated abrasive (like a sandpaper) stuck to a soft, non-slip surface to make achieving a convex grind easier. Similarly, a brass mirror may be cut with a bonded abrasive, have its surface flattened with a coated abrasive to achieve a basic shape, and then have finer grades of abrasive successively applied culminating in a wax paste impregnated with rouge to leave a sort of "grainless finish" called, in this case, a "mirror finish".
Also, different shapes of adhesive may make it harder to abrade certain areas of the workpiece. Health hazards can arise from any dust produced (which may be ameliorated through the use of a lubricant) which could lead to silicosis (when the abrasive or workpiece is a silicate) and the choice of any lubricant. Besides water, oils are the most common lubricants. These may present inhalation hazards, contact hazards and, as friction necessarily produces heat, flammable material hazards.
An abrasive which is too hard or too coarse can remove too much material or leave undesired scratch marks. Besides being unsightly, scratching can have other, more serious effects. Excessive abrasion or the presence of scratches may:
diminish or destroy usefulness (as in the case of scratching optical lenses and compact discs or dulling knives);
trap dirt, water, or other material;
increase surface area (permitting greater chemical reactivity such as increased rusting which is also affected by matter caught in scratches);
erode or penetrate a coating (such as a paint or a chemical or wear resistant coating);
overly quickly cause an object to wear away (such as a blade or a gemstone);
increase friction (as in jeweled bearings and pistons).
A finer or softer abrasive will tend to leave much finer scratch marks which may even be invisible to the naked eye (a "grainless finish"); a softer abrasive may not even significantly abrade a certain object. A softer or finer abrasive will take longer to cut, as it tends to cut less deeply than a coarser, harder material. Also, the softer abrasive may become less effective more quickly as the abrasive is itself abraded. This allows fine abrasives to be used in the polishing of metal and lenses where the series of increasingly fine scratches tends to take on a much more shiny or reflective appearance or greater transparency. Very fine abrasives may be used to coat the strop for a cut-throat razors, however, the purpose of stropping is not to abrade material but to straighten the burr on an edge. The final stage of sharpening Japanese swords is called polishing and may be a form of superfinishing.
Different chemical or structural modifications may be made to alter the cutting properties of the abrasive.
Other very important considerations are price and availability. Diamond, for a long time considered the hardest substance in existence, is actually softer than fullerite and even harder aggregated diamond nanorods, both of which have been synthesised in laboratories, but no commercial process has yet been developed. Diamond itself is expensive due to scarcity in nature and the cost of synthesising it. Bauxite is a very common ore which, along with corundum's reasonably high hardness, contributes to corundum's status as a common, inexpensive abrasive.
Thought must be given to the desired task about using an appropriately hard abrasive. At one end, using an excessively hard abrasive wastes money by wearing it down when a cheaper, less hard abrasive would suffice. At the other end, if the abrasive substance is too soft, abrasion does not take place in a timely fashion, effectively wasting the abrasive as well as any accruing costs associated with loss of time.
Other instances of abrasion
Aside from the aforementioned uses of shaping and finishing, abrasives may also be used to prepare surfaces for application of some sort of paint of adhesive. An excessively smooth surface may prevent paint and adhesives from adhering as strongly as an irregular surface could allow. Inflatable tyre repair kits (which, on bicycles particularly, are actually patches for the inner tube rather than the tyre) require use of an abrasive so that the self-vulcanising cement will stick strongly.
Inadvertently, people who use knives on glass or metal cutting boards are abrading their knife blades. The pressure at the knife edge can easily create microscopic (or even macroscopic) cuts in the board. This cut is a ready source of abrasive material as well as a channel full of this abrasive through which the edge slides. For this reason, and without regard for the health benefits, wooden boards are much more desirable. A similar occurrence arises with glass-cutters. Glass-cutters have circular blades that are designed to roll not slide. They should never retrace an already effected cut.
Undesired abrasion may result from the presence of carbon in internal combustion engines. While smaller particles are readily transported by the lubrication system, larger carbon particles may abrade components with close tolerances. The carbon arises from the excessive heating of engine oil or from incomplete combustion. This soot may contain fullerenes which are noted for their extreme hardness—and small size and limited quantity which would tend to limit their effect.
| Technology | Material and chemical | null |
346547 | https://en.wikipedia.org/wiki/Reusability | Reusability | In computer science and software engineering, reusability is the use of existing assets in some form within the software product development process; these assets are products and by-products of the software development life cycle and include code, software components, test suites, designs and documentation. The opposite concept of reusability is leverage, which modifies existing assets as needed to meet specific system requirements. Because reuse implies the creation of a , it is preferred over leverage.
Subroutines or functions are the simplest form of reuse. A chunk of code is regularly organized using modules or namespaces into layers. Proponents claim that objects and software components offer a more advanced form of reusability, although it has been tough to objectively measure and define levels or scores of reusability.
The ability to reuse relies in an essential way on the ability to build larger things from smaller parts, and being able to identify commonality among those parts. Reusability is often a required characteristic of platform software. Reusability brings several aspects to software development that do not need to be considered when reusability is not required.
Reusability implies some explicit management of build, packaging, distribution, installation, configuration, deployment, maintenance and upgrade issues. If these issues are not considered, software may appear to be reusable from design point of view, but will not be reused in practice.
Software reusability more specifically refers to design features of a software element (or collection of software elements) that enhance its suitability for reuse.
Many reuse design principles were developed at the WISR workshops.
Candidate design features for software reuse include:
Adaptable
Brief: small size
Consistency
Correctness
Extensibility
Fast
Flexible
Generic
Localization of volatile (changeable) design assumptions (David Parnas)
Modularity
Orthogonality
Simple: low complexity
Stability under changing requirements
Consensus has not yet been reached on this list on the relative importance of the entries nor on the issues which make each one important for a particular class of applications.
| Technology | Software development: General | null |
346860 | https://en.wikipedia.org/wiki/Ajwain | Ajwain | Ajwain or ajowan (Trachyspermum ammi) () —also known as ajowancaraway,
thymol seeds, bishop's weed, or carom—is an annual herb in the family Apiaceae. Both the leaves and the seed‑like fruit (often mistakenly called seeds) of the plant are consumed by humans. The name "bishop's weed" also is a common name for other plants. The "seed" (i.e., the fruit) is often confused with lovage seed.
Description
Ajwain's small, oval, seed-like fruits are pale brown schizocarps, which resemble the seeds of other plants in the family Apiaceae such as caraway, cumin and fennel. They have a bitter and pungent taste, with a flavor similar to anise and oregano. They smell like thyme because they also contain thymol, but they are more aromatic and less subtle in taste, as well as being somewhat bitter and pungent. Even a small number of fruits tend to dominate the flavor of a dish.
Cultivation and production
Ajwain grows in dry, barren soil in its indigenous regions of India, Iran, Afghanistan, and parts of northern Africa. Gujarat and Rajasthan are regions in India well-known for cultivating ajwain.
Culinary uses
The fruits are rarely eaten raw; they are commonly dry-roasted or fried in ghee (clarified butter). This allows the spice to develop a more subtle and complex aroma. It is widely used in the cuisine of the Indian subcontinent, often as part of a chaunk (also called a tarka), a mixture of spices – sometimes with a little chopped garlic or onion – fried in oil or clarified butter, which is used to flavor a dish at the end of cooking. In Afghanistan, the fruits are sprinkled over bread and biscuits.
Other applications of ajwain include incorporating the seeds in specific types of breads, such as naans and parathas. The seeds can also be used as a mouth freshener when mixed with lemon juice and black pepper, and then dried, or can be used as an ingredient in hot tea.
In herbalism
Ajwain is used in herbalism practices, such as Ayurveda, in the belief that it can treat various disorders. However, there is no good evidence that ajwain is effective as a therapy for treating any disease.
Adverse effects
Pregnant women should avoid ajwain due to potential adverse effects on fetal development, and its use is discouraged while breastfeeding. In high amounts taken orally, ajwain can result in fatal poisoning. People taking nonsteroidal anti-inflammatory drugs or antiplatelet medications are susceptible to adverse effects from ajwain ingestion, as ajwain has anti-clotting activity of its own.
Essential oil
Hydrodistillation of ajwain fruits yields an essential oil consisting primarily of thymol, gamma-terpinene, p-cymene, and more than 20 trace compounds which are predominantly terpenoids.
| Biology and health sciences | Herbs and spices | Plants |
346865 | https://en.wikipedia.org/wiki/Basil | Basil | Basil (, ; , ; Ocimum basilicum (, )), also called great basil, is a culinary herb of the family Lamiaceae (mints). It is a tender plant, and is used in cuisines worldwide. In Western cuisine, the generic term "basil" refers to the variety also known as Genovese basil or sweet basil. Basil is native to tropical regions from Central Africa to Southeast Asia. In temperate climates basil is treated as an annual plant, but it can be grown as a short-lived perennial or biennial in warmer horticultural zones with tropical or Mediterranean climates.
There are many varieties of basil including sweet basil, Thai basil (O. basilicum var. thyrsiflora), and Mrs. Burns' Lemon (O.basilicum var. citriodora). O. basilicum can cross-pollinate with other species of the Ocimum genus, producing hybrids such as lemon basil (O. × citriodorum) and African blue basil (O. × kilimandscharicum).
Description
Basil is an annual, or sometimes perennial, herb. Depending on the variety, plants can reach heights of between . Basil leaves are glossy and ovulate, with smooth or slightly toothed edges that typically cup slightly; the leaves are arranged oppositely along the square stems. Leaves may be green or purple. Its flowers are small and white, and grow from a central inflorescence, or spike, that emerges from the central stem atop the plant. Unusual among Lamiaceae, the four stamens and the pistil are not pushed under the upper lip of the corolla, but lie over the inferior lip. After entomophilous pollination, the corolla falls off and four round achenes develop inside the bilabiate calyx.
Phytochemistry
The various basils have such distinct scents because the volatile aromatic compounds vary with cultivars. The essential oil from European basil contains high concentrations of linalool and methyl chavicol (estragole), in a ratio of about 3:1. Other constituents include: 1,8-cineole, eugenol, and myrcene, among others. The clove scent of sweet basil is derived from eugenol. The aroma profile of basil includes 1,8-cineole and methyl eugenol. In this species eugenol is synthesised from coniferyl acetate and NADPH.
Similar species
Some similar species in the same genus may be commonly called "basil", although they are not varieties of Ocimum basilicum.
Camphor basil, African basil (O. kilimandscharicum)
Clove basil, also African basil (Ocimum gratissimum)
Holy basil (Ocimum tenuiflorum, formerly known as O. sanctum)
Taxonomy
The exact taxonomy of basil is uncertain due to the immense number of cultivars, its ready polymorphy, and frequent cross-pollination (resulting in new hybrids) with other members of the genus Ocimum and within the species. Ocimum basilicum has at least 60 varieties, which further complicates taxonomy.
Cultivars
Most basils are cultivars of sweet basil. Most basil varieties have green leaves, but a few are purple, such as, 'Purple Delight'.
Anise basil, Licorice basil, or Persian basil (O. basilicum 'Liquorice')
Cinnamon basil (Ocimum basilicum 'Cinnamon')
Dark opal basil (Ocimum basilicum 'Dark Opal')
Genovese basil or Sweet Basil (Ocimum basilicum)
Greek basil (Ocimum basilicum var. minimum)
Globe basil, dwarf basil, French basil (Ocimum basilicum 'Minimum')
Lettuce leaf basil (Ocimum basilicum 'Crispum')
Napolitano basil, also known as Napoletano basil, Neapolitan basil, Mammoth basil, Bolloso Napoletano basil, Napolitano Mammoth-Leafed basil, or Italian Large-Leaf basil (Ocimum basilicum)
Purple basil (Ocimum basilicum 'Purpurescens')
Rubin basil (Ocimum basilicum 'Rubin')
Thai basil (Ocimum basilicum thyrsifolium)
Hybrids
African blue basil (Ocimum basilicum × O. kilimandscharicum)
Lemon basil (Ocimum basilicum × O. americanum)
Spice basil (Ocimum basilicum × O. americanum), which is sometimes sold as holy basil
Etymology
The name "basil" comes from the Latin , and the Greek (), meaning "royal/kingly plant", possibly because the plant was believed to have been used in production of royal perfumes. Basil is likewise sometimes referred to in French as ('the royal herb'). The Latin name has been confused with basilisk, as it was supposed to be an antidote to the basilisk's venom.
Distribution and habitat
Basil is native to India and other tropical regions stretching from Africa to South East Asia, but has now become globalized due to human cultivation.
Cultivation
Growing conditions
Basil is sensitive to cold, with best growth in hot, dry conditions. It behaves as an annual if there is any chance of a frost. However, due to its popularity, basil is cultivated in many countries around the world. Production areas include countries in the Mediterranean area, those in the temperate zone, and others in subtropical climates.
In Northern Europe, Canada, the northern states of the U.S., and the South Island of New Zealand, basil grows best if sown under glass in a peat pot, then planted out in late spring/early summer (when there is little chance of a frost); however, it can also thrive when planted outside in these climates. Additionally, it may be sown in soil once chance of frost is past. It fares best in well-drained soil with direct exposure to the sun.
Although basil grows best outdoors, it can be grown indoors in a pot and, like most herbs, will do best on a sun-facing windowsill, kept away from extremely cold drafts. A greenhouse or row cover is ideal if available. It can, however, even be grown in a basement under fluorescent lights. Supplemental lighting produces greater biomass and phenol production, with red + blue specifically increasing growth and flower bud production. increases the volatiles in O. basilicum essential oil, which has not been reproducible in other plants, and so may be unique to the genus or even to this species.
Basil plants require regular watering, but not as much attention as is needed in other climates. If its leaves have wilted from lack of water, it will recover if watered thoroughly and placed in a sunny location. Yellow leaves towards the bottom of the plant are an indication that the plant has been stressed; usually this means that it needs less water, or less or more fertilizer. Basil can be propagated reliably from cuttings with the stems of short cuttings suspended in water for two weeks or until roots develop.
Pruning, flowering, and seeding
Once a stem produces flowers, foliage production stops on that stem, the stem becomes woody, and essential oil production declines. To prevent this, a basil-grower may pinch off any flower stems before they are fully mature. Because only the blooming stem is so affected, some stems can be pinched for leaf production, while others are left to bloom for decoration or seeds. Picking the leaves off the plant helps promote growth, largely because the plant responds by converting pairs of leaflets next to the topmost leaves into new stems.
Once the plant is allowed to flower, it may produce seed pods containing small black seeds, which can be saved and planted the following year. If allowed to go to seed, a basil plant will grow back the next year.
Diseases
Basil suffers from several plant pathogens that can ruin the crop and reduce yield. Fusarium wilt is a soil-borne fungal disease that will quickly kill younger basil plants. Seedlings may be killed by Pythium damping off. A common foliar disease of basil is gray mold caused by Botrytis cinerea; it can cause infections post-harvest and is capable of killing the entire plant. Black spot can be seen on basil foliage and is caused by the fungi genus Colletotrichum. Downy mildew caused by Peronospora belbahrii is a significant disease, as first reported in Italy in 2003. It was reported in the Florida in 2007 and by 2008 had spread along the eastern United States, reaching Canada. Basil cultivars resistant to P. belbahrii have been developed.
Non-pathogenic bacteria found on basil include Novosphingobium species.
Uses
Culinary
Basil is most commonly used fresh in recipes. In general, it is added last, as cooking quickly destroys the flavor. The fresh herb can be kept for a short time in plastic bags in the refrigerator, or for a longer period in the freezer, after being blanched quickly in boiling water.
Leaves and flowers
The most commonly used Mediterranean basil cultivars are "Genovese", "Purple Ruffles", "Mammoth", "Cinnamon", "Lemon", "Globe", and "African Blue". Basil is one of the main ingredients in pesto, an Italian sauce with olive oil and basil as its primary ingredients. Many national cuisines use fresh or dried basils in soups and other foods, such as to thicken soups. Basil is commonly steeped in cream or milk to create flavor in ice cream or chocolate truffles.
Lemon basil has a strong lemony smell and flavor due to the presence of citral. It is widely used in Indonesia, where it is called and served raw as an accompaniment to meat or fish.
Seeds
When soaked in water, the seeds of several basil varieties become gelatinous, and are used in Asian drinks and desserts such as the Indian faluda, the Iranian , or . In Kashmir, the Ramadan fast is often broken with babre beole, a sharbat made with basil seeds.
Folk medicine
Basil is used in folk medicine practices, such as those of Ayurveda or traditional Chinese medicine.
Insecticide and insect repellent
Studies of the essential oil have shown insecticidal and insect-repelling properties, including potential toxicity to mosquitos. The essential oil is found by Huignard et al. 2008 to inhibit electrical activity by decreasing action potential amplitude, by shortening the post hyperpolarization phase, and reducing the action frequency of action potentials. In Huignard's opinion this is due to the linalool and estragole, the amplitude reduction due to linalool, and the phase shortening due to both.
Callosobruchus maculatus, a pest which affects cowpea, is repelled by the essential oil. The essential oil mixed with kaolin is both an adulticide and an ovicide, effective for three months against C. maculatus in cowpea. The thrips Frankliniella occidentalis and Thrips tabaci are repelled by O. basilicum, making this useful as an insect repellent in other crops. The pests Sitophilus oryzae, Stegobium paniceum, Tribolium castaneum, and Bruchus chinensis are evaluated by Deshpande et al. 1974 and '77.
Nematicide
The essential oil is found by Malik et al. 1987 and Sangwan et al. 1990 to be nematicidal against Tylenchulus semipenetrans, Meloidogyne javanica, Anguina tritici, and Heterodera cajani.
Bacterial and fungal inhibition
The essential oil of the leaf and terminal shoot is effective against a large number of bacterial species including Lactiplantibacillus plantarum and Pseudomonas spp. The essential oil of the leaf and terminal shoot is also effective against a large number of fungal species including Aspergillus spp., Candida spp., Mucor spp., and Geotrichum candidum.
In culture
Religion
There are many rituals and beliefs associated with basil. The ancient Egyptians and ancient Greeks believed basil would open the gates of heaven for a person passing on. However, Herbalist Nicholas Culpeper saw basil as a plant of dread and suspicion.
In Portugal, dwarf bush basil is traditionally presented in a pot, together with a poem and a paper carnation, to a sweetheart, on the religious holidays of John the Baptist (see ) and Saint Anthony of Padua.
Basil has religious significance in the Greek Orthodox Church, where it is used to sprinkle holy water. The Bulgarian Orthodox Church, Serbian Orthodox Church, Macedonian Orthodox Church and Romanian Orthodox Church use basil (, ; , ; , ) to prepare holy water and pots of basil are often placed below church altars. Some Greek Orthodox Christians avoid eating it due to its association with the legend of the Elevation of the Holy Cross.
Art and literature
In Giovanni Boccaccio's 14th century Decameron, the fifth story of the narrative's fourth day involves a pot of basil as a central plot device. This famous story inspired John Keats to write his 1814 poem "Isabella, or the Pot of Basil", which was in turn the inspiration for two paintings of the Pre-Raphaelite Brotherhood: John Everett Millais's Isabella in 1849 and in 1868 the Isabella and the Pot of Basil by William Holman Hunt.
| Biology and health sciences | Lamiales | null |
346892 | https://en.wikipedia.org/wiki/Tendril | Tendril | In botany, a tendril is a specialized stem, leaf or petiole with a thread-like shape used by climbing plants for support and attachment, as well as cellular invasion by parasitic plants such as Cuscuta. There are many plants that have tendrils; including sweet peas, passionflower, grapes and the Chilean glory-flower. Tendrils respond to touch and to chemical factors by curling, twining, or adhering to suitable structures or hosts. Tendrils vary greatly in size from a few centimeters up to 27 inches (69 centimeters) for Nepenthes harryana The chestnut vine (Tetrastigma voinierianum) can have tendrils up to 20.5 inches (52 centimeters) in length. Normally there is only one simple or branched tendril at each node (see plant stem), but the aardvark cucumber (Cucumis humifructus) can have as many as eight.
History
The earliest and most comprehensive study of tendrils was Charles Darwin's monograph On the Movements and Habits of Climbing Plants, which was originally published in 1865. This work also coined the term circumnutation to describe the motion of growing stems and tendrils seeking supports. Darwin also observed the phenomenon now known as tendril perversion, in which tendrils adopt the shape of two sections of counter-twisted helices with a transition in the middle.
Biology of tendrils
In the garden pea, it is only the terminal leaflets that are modified to become tendrils. In other plants such as the yellow vetch (Lathyrus aphaca), the whole leaf is modified to become tendrils while the stipules become enlarged and carry out photosynthesis. Still others use the rachis of a compound leaf as a tendril, such as members of the genus Clematis.
The specialised pitcher traps of Nepenthes plants form on the end of tendrils. The tendrils of aerial pitchers are usually coiled in the middle. If the tendril comes into contact with an object for long enough it will usually curl around it, forming a strong anchor point for the pitcher. In this way, the tendrils help to support the growing stem of the plant. Tendrils of Cuscuta, a parasitic plant, are guided by airborne chemicals, and only twine around suitable hosts.
Evolution and species
Climbing habits in plants support themselves to reach the canopy in order to receive more sunlight resources and increase the diversification in flowering plants. Tendril is a plant organ that is derived from various morphological structures such as stems, leaves and inflorescences. Even though climbing habits are involved in the angiosperms, gymnosperms, and fern, tendrils are often shown in angiosperms and little in fern. Based on their molecular basis of tendril development, studies showed that tendrils helical growth performance is not correlated with ontogenetic origin, instead, there are multiple ontogenetic origins. 17 types of tendrils have been identified by their ontogenetic origins and growth pattern, and each type of tendril can be involved more than once within angiosperms. Common fruits and vegetables that have tendrils includes watermelon (Citrullus lanatus)'s derived from modified stem, pea (Pisum sativum)'s derived from modified terminal leaflets and common grape vine (Vitis vinifera)'s is modified from whole inflorescence.
Coiling mechanism
Circumnutation
The mechanism of tendril coiling begins with circumnutation of the tendril in which it is moving and growing in a circular oscillatory pattern around its axis. Circumnutation is often defined as the first main movement of the tendril, and it serves the purpose of increasing the chance that the plant will come in contact with a support system (physical structure for the tendril to coil around). In a 2019 study done by Guerra et al., it was shown that without a support stimulus, in this case a stake in the ground, the tendrils will circumnutate towards a light stimulus. After many attempts to reach a support structure, the tendril will eventually fall to the ground. However, it was found that when a support stimulus is present, the tendril’s circumnutation oscillation occurs in the direction of the support stimulus. Therefore, it was concluded that tendrils are able to change the direction of their circumnutation based on the presence of a support stimulus. The process of circumnutation in plants is not unique to tendril plants, as almost all plant species show circumnutation behaviors.
Contact coiling
Thigmotropism is the basis of the input signal in the tendril coiling mechanism. For example, pea tendrils have highly sensitive cells in the surfaces of cell walls that are exposed. These sensitized cells are the ones that initiate the thigmotropic signal, typically as a calcium wave. The primary touch signal induces a signaling cascade of other phytohormones, most notably gamma-Aminobutyric acid (GABA) and Jasmonate (JA). In grapevine tendrils, it recently has been shown that GABA can independently promote tendril coiling. It has also been shown that jasmonate phytohormones serve as a hormonal signal to initiate tendril coiling. This cascade can activate plasma membrane H+-ATPase, which also plays a role in the contact coiling mechanism as a proton pump. This pump activity establishes an electrochemical of H+ ions from inside the cell to the apoplast, which in turn creates an osmotic gradient. This leads to loss of turgor pressure; the differences in cell size due to the loss of turgor pressure in some cells creates the coiling response. This contractile movement is also influenced by gelatinous fibers, which contract and lignify in response to the thigmotropic signal cascade.
Self-discrimination
Although tendrils twine around hosts based on touch perception, plants have a form of self-discrimination and avoid twining around themselves or neighboring plants of the same speciesdemonstrating chemotropism based on chemoreception. Once a tendril comes in contact with a neighboring conspecific plant (of the same species) signaling molecules released by the host plant bind to chemoreceptors on the climbing plant’s tendrils. This generates a signal that prevents the thigmotropic pathway and therefore prevents the tendril from coiling around that host.
Studies confirming this pathway have been performed on the climbing plant Cayratia japonica. Research demonstrated that when two C. japonica plants were placed in physical contact, the tendrils would not coil around the conspecific plant. Researchers tested this interaction by isolating oxalate crystals from the leaves of a C. japonica plant and coating a stick with the oxalate crystals. The tendrils of C. japonica plants that came in physical contact with the oxalate-coated stick would not coil, confirming that climbing plants use chemoreception for self-discrimination.
Self-discrimination may confer an evolutionary advantage for climbing plants to avoid coiling around conspecific plants. This is because neighboring climbing plants do not provide as stable of structures to coil around when compared to more rigid nearby plants. Furthermore, by being able to recognize and avoid coiling around conspecific plants, the plants reduce their proximity to competition, allowing them to have access to more resources and therefore better growth.
Gallery
| Biology and health sciences | Plant anatomy and morphology: General | Biology |
346896 | https://en.wikipedia.org/wiki/Berkeley%20Open%20Infrastructure%20for%20Network%20Computing | Berkeley Open Infrastructure for Network Computing | The Berkeley Open Infrastructure for Network Computing (BOINC, pronounced – rhymes with "oink") is an open-source middleware system for volunteer computing (a type of distributed computing). Developed originally to support SETI@home, it became the platform for many other applications in areas as diverse as medicine, molecular biology, mathematics, linguistics, climatology, environmental science, and astrophysics, among others. The purpose of BOINC is to enable researchers to utilize processing resources of personal computers and other devices around the world.
BOINC development began with a group based at the Space Sciences Laboratory (SSL) at the University of California, Berkeley, and led by David P. Anderson, who also led SETI@home. As a high-performance volunteer computing platform, BOINC brings together 34,236 active participants employing 136,341 active computers (hosts) worldwide, processing daily on average 20.164 PetaFLOPS (it would be the 21st largest processing capability in the world compared with an individual supercomputer). The National Science Foundation (NSF) funds BOINC through awards SCI/0221529, SCI/0438443 and SCI/0721124. Guinness World Records ranks BOINC as the largest computing grid in the world.
BOINC code runs on various operating systems, including Microsoft Windows, macOS, Android, Linux, and FreeBSD. BOINC is free software released under the terms of the GNU Lesser General Public License (LGPL).
History
BOINC was originally developed to manage the SETI@home project. David P. Anderson has said that he chose its name because he wanted something that was not "imposing", but rather "light, catchy, and maybe - like 'Unix' - a little risqué", so he "played around with various acronyms and settled on 'BOINC'".
The original SETI client was a non-BOINC software exclusively for SETI@home. It was one of the first volunteer computing projects, and not designed with a high level of security. As a result, some participants in the project attempted to cheat the project to gain "credits", while others submitted entirely falsified work. BOINC was designed, in part, to combat these security breaches.
The BOINC project started in February 2002, and its first version was released on April 10, 2002. The first BOINC-based project was Predictor@home, launched on June 9, 2004. In 2009, AQUA@home deployed multi-threaded CPU applications for the first time, followed by the first OpenCL application in 2010.
As of 15 August 2022, there are 33 projects on the official list. There are also, however, BOINC projects not included on the official list. Each year, an international BOINC Workshop is hosted to increase collaboration among project administrators. In 2021, the workshop was hosted virtually.
While not affiliated with BOINC officially, there have been several independent projects that reward BOINC users for their participation, including Charity Engine (sweepstakes based on processing power with prizes funded by private entities who purchase computational time of CE users), Bitcoin Utopia (now defunct), and Gridcoin (a blockchain which mints coins based on processing power).
Design and structure
BOINC is software that can exploit the unused CPU and GPU cycles on computer hardware to perform scientific computing. In 2008, BOINC's website announced that Nvidia had developed a language called CUDA that uses GPUs for scientific computing. With NVIDIA's assistance, several BOINC-based projects (e.g., MilkyWay@home. SETI@home) developed applications that run on NVIDIA GPUs using CUDA. BOINC added support for the ATI/AMD family of GPUs in October 2009. The GPU applications run from 2 to 10 times faster than the former CPU-only versions. GPU support (via OpenCL) was added for computers using macOS with AMD Radeon graphic cards, with the current BOINC client supporting OpenCL on Windows, Linux, and macOS. GPU support is also provided for Intel GPUs.
BOINC consists of a server system and client software that communicate to process and distribute work units and return results.
Mobile application
A BOINC app also exists for Android, allowing every person owning an Android device – smartphone, tablet and/or Kindle – to share their unused computing power. The user is allowed to select the research projects they want to support, if it is in the app's available project list.
By default, the application will allow computing only when the device is connected to a WiFi network, is being charged, and the battery has a charge of at least 90%. Some of these settings can be changed to users needs. Not all BOINC projects are available and some of the projects are not compatible with all versions of Android operating system or availability of work is intermittent. Currently available projects are Asteroids@home, Einstein@Home, LHC@home, Moo! Wrapper, Rosetta@home, World Community Grid and . As of September 2021, the most recent version of the mobile application can only be downloaded from the BOINC website or the F-Droid repository as the official Google Play store does not allow downloading and running executables not signed by the app developer and each BOINC project has their own executable files.
User interfaces
BOINC can be controlled remotely by remote procedure calls (RPC), from the command line, and from a BOINC Manager. BOINC Manager currently has two "views": the Advanced View and the Simplified GUI. The Grid View was removed in the 6.6.x clients as it was redundant. The appearance (skin) of the Simplified GUI is user-customizable, in that users can create their own designs.
Account managers
A BOINC Account Manager is an application that manages multiple BOINC project accounts across multiple computers (CPUs) and operating systems. Account managers were designed for people who are new to BOINC or have several computers participating in several projects. The account manager concept was conceived and developed jointly by GridRepublic and BOINC. Current and past account managers include:
BAM! (BOINC Account Manager) (The first publicly available Account Manager, released for public use on May 30, 2006)
GridRepublic (Follows the ideas of simplicity and neatness in account management)
Charity Engine (Non-profit account manager for hire, uses prize draws and continuous charity fundraising to motivate people to join the grid)
Science United (An account manager designed to make BOINC easier to use which automatically selects vetted BOINC projects for users based on desired research areas such as "medicine" or "physics")
Dazzler (Open-source Account Manager, to ease institutional management resources)
Credit system
The BOINC Credit System is designed to avoid bad hardware and cheating by validating results before granting credit.
The credit management system helps to ensure that users are returning results which are both statistically and scientifically accurate.
Online volunteer computing is a complicated and variable mix of long-term users, retiring users and new users with different personal aspirations.
Projects
BOINC is used by many groups and individuals. Some BOINC projects are based at universities and research labs while others are independent areas of research or interest.
Active
Completed
| Technology | Science | null |
346939 | https://en.wikipedia.org/wiki/Carob | Carob | The carob ( ; Ceratonia siliqua) is a flowering evergreen tree or shrub in the Caesalpinioideae sub-family of the legume family, Fabaceae. It is widely cultivated for its edible fruit, which takes the form of seed pods, and as an ornamental tree in gardens and landscapes. The carob tree is native to the Mediterranean region and the Middle East. Portugal is the largest producer of carob, followed by Italy and Morocco.
In the Mediterranean Basin, extended to the southern Atlantic coast of Portugal (i.e., the Algarve region) and the Atlantic northwestern Moroccan coast, carob pods were often used as animal feed and in times of famine, as "the last source of [human] food in hard times". The ripe, dried, and sometimes toasted pod is often ground into carob powder, which was sometimes used as a substitute for cocoa powder, especially in the 1970s natural food movement. The powder and chips can be used as a chocolate alternative in most recipes.
The plant's seeds are used to produce locust bean gum or carob gum, a common thickening agent used in food processing.
Description
The carob tree grows up to tall. The crown is broad and semispherical, supported by a thick trunk with rough brown bark and sturdy branches. Its leaves are long, alternate, pinnate, and may or may not have a terminal leaflet. It is frost-tolerant to roughly .
Most carob trees are dioecious and some are hermaphroditic, so strictly male trees do not produce fruit. When the trees blossom in autumn, the flowers are small and numerous, spirally arranged along the inflorescence axis in catkin-like racemes borne on spurs from old wood and even on the trunk (cauliflory); they are pollinated by both wind and insects. The male flowers smell like human semen, an odor that is caused in part by amines.
The fruit is a legume (also known commonly, but less accurately, as a pod), that is elongated, compressed, straight, or curved, and thickened at the sutures. The pods take a full year to develop and ripen. When the sweet, ripe pods eventually fall to the ground, they are eaten by various mammals, such as swine, thereby dispersing the hard inner seed in the excrement.
The seeds of the carob tree contain leucodelphinidin, a colourless flavanol precursor related to leucoanthocyanidins.
Etymology
The word "carob" comes from Middle French (modern French ), which borrowed it from Arabic (kharrūb, "locust bean pod") and Persian khirnub, which ultimately borrowed it perhaps from Akkadian language or Aramaic ḥarrūḇā. '
Ceratonia siliqua, the scientific name of the carob tree, derives from the Greek keratōnia, "carob-tree" (cf. kéras, "horn"), and Latin siliqua "pod, carob".
In English, it is also known as "St. John's bread" and "locust tree" (not to be confused with African locust bean). The latter designation also applies to several other trees from the same family.
In Yiddish, it is called bokser, derived from the Middle High German bokshornboum "ram's horn tree" (in reference to the shape of the carob).
The carat, a unit of mass for gemstones, and a measurement of purity for gold, takes its name via the Arabic qīrāṭ from the Greek name for the carob seed (lit. "small horn").
Distribution and habitat
Although cultivated extensively, carob can still be found growing wild in eastern Mediterranean regions, and has become naturalized in the western Mediterranean.
The tree is typical in the southern Portuguese region of the Algarve, where the tree is called alfarrobeira, and the fruit alfarroba. It is also seen in southern and eastern Spain (, Catalan / Valencian / Balearic: garrofer, garrofera, garrover, garrovera), mainly in the regions of Andalusia, Murcia, Valencia, the Balearic Islands and Catalonia (Catalan / Valencian / Balearic: garrofer, garrofera, garrover, garrovera); Malta (), on the Italian islands of Sicily () and Sardinia (), in Southern Croatia (), such as on the island of Šipan, in eastern Bulgaria (), and in Southern Greece, Cyprus, as well as on many Greek islands such as Crete and Samos.
In Israel, the Hebrew name is חרוב (translit. charuv). The common Greek name is (translit. ), or (translit. , meaning "wooden horn"). In Turkey, it is known as "goat's horn" ().
The various trees known as algarrobo in Latin America (Samanea saman in Cuba, Prosopis pallida in Peru, and four species of Prosopis in Argentina and Paraguay) belong to a different subfamily of the Fabaceae: Mimosoideae. Early Spanish settlers named them algarrobo after the carob tree because they also produce pods with sweet pulp.
Ecology
The carob genus, Ceratonia, belongs to the legume family, Fabaceae, and is believed to be an archaic remnant of a part of this family now generally considered extinct. It grows well in warm temperate and subtropical areas, and tolerates hot and humid coastal areas. As a xerophyte (drought-resistant species), carob is well adapted to the conditions of the Mediterranean region with just of rainfall per year.
Carob trees can survive long periods of drought, but to grow fruit, they need of rainfall per year. They prefer well-drained, sandy loams and are intolerant of waterlogging, but the deep root systems can adapt to a wide variety of soil conditions and are fairly salt-tolerant (up to 3% in soil). After being irrigated with saline water in the summer, carob trees could possibly recover during winter rainfalls. In some experiments, young carob trees were capable of basic physiological functions under high-salt conditions (40 mmol NaCl/L).
Not all legume species can develop a symbiotic relationship with rhizobia to make use of atmospheric nitrogen. It remains unclear if carob trees have this ability: Some findings suggest that it is not able to form root nodules with rhizobia, while in another more recent study, trees have been identified with nodules containing bacteria believed to be from the genus Rhizobium. However, a study measuring the 15N-signal (isotopic signature) in the tissue of the carob tree did not support the theory that carob trees naturally use atmospheric nitrogen.
Cultivation
The vegetative propagation of carob is naturally restricted due to its low adventitious rooting potential. Therefore, grafting and air-layering may prove to be more effective methods of asexual propagation. Seeds are commonly used as the propagation medium. The sowing occurs in pot nurseries in early spring and the cooling- and drying-sensitive seedlings are then transplanted to the field in the next year after the last frost. Carob trees enter slowly into production phase. Where in areas with favorable growing conditions, the cropping starts 3–4 years after budding, with the nonbearing period requiring up to 8 years in regions with marginal soils. Full bearing of the trees occurs mostly at a tree-age of 20–25 years when the yield stabilizes. The orchards are traditionally planted in low densities of 25–45 trees per hectare (). Hermaphroditic or male trees, which produce fewer or no pods, respectively, are usually planted in lower densities in the orchards as pollenizers.
Intercropping with other tree species is widely spread. Not much cultivation management is required. Only light pruning and occasional tilling to reduce weeds is necessary. Nitrogen-fertilizing of the plants has been shown to have positive impacts on yield performance. Although it is native to moderately dry climates, two or three summers' irrigation greatly aid the development, hasten the fruiting, and increase the yield of a carob tree.
Harvest and post-harvest treatment
The most labour-intensive part of carob cultivation is harvesting, which is often done by knocking the fruit down with a long stick and gathering them together with the help of laid-out nets. This is a delicate task because the trees are flowering at the same time and care has to be taken not to damage the flowers and the next year's crop. The literature recommends research to get the fruit to ripen more uniformly or also for cultivars which can be mechanically harvested (by shaking).
After harvest, carob pods have a moisture content of 10–20% and should be dried down to a moisture content of 8% so the pods do not rot. Further processing separates the kernels (seeds) from the pulp. This process is called kibbling and results in seeds and pieces of carob pods (kibbles). Processing of the pulp includes grinding for animal feed production or roasting and milling for human food industry. The seeds have to be peeled which happens with acid or through roasting. Then the endosperm and the embryo are separated for different uses.
Pests and diseases
Few pests are known to cause severe damage in carob orchards, so they have traditionally not been treated with pesticides. Some generalist pests such as the larvae of the leopard moth (Zeuzera pyrina L.), the dried fruit moth (Cadra calidella), small rodents such as rats (Rattus spp.) and gophers (Pitymys spp.) can cause damage occasionally in some regions. Only some cultivars are severely susceptible to mildew disease (Oidium ceratoniae C.). One pest directly associated with carob is the larva of the carob moth (Myelois ceratoniae Z.), which can cause extensive postharvest damage.
Cadra calidella attack carob crops before harvest and infest products in stores. This moth, prevalent in Cyprus, will often infest the country's carob stores. Research has been conducted to understand the physiology of the moth, in order to gain insight on how to monitor moth reproduction and lower their survival rates, such as through temperature control, pheromone traps, or parasitoid traps.
Production
In 2022, world production of carob (as locust beans) was estimated to be 56,423 tonnes, although not all countries known to grow carob reported their results to the UN Food and Agriculture Organization. Production amounts for Turkey and Morocco accounted for nearly all the world total reported in 2022.
Cultivars and breeding aims
Most of the roughly 50 known cultivars are of unknown origin and only regionally distributed. The cultivars show high genetic and therefore morphological and agronomical variation. No conventional breeding by controlled crossing has been reported, but selection from orchards or wild populations has been done. Domesticated carobs (C. s. var. edulis) can be distinguished from their wild relatives (C. s. var. silvestris) by some fruit-yielding traits such as building of greater beans, more pulp, and higher sugar contents. Also, genetic adaptation of some varieties to the climatic requirements of their growing regions has occurred. Though a partially successful breaking of the dioecy happened, the yield of hermaphrodite trees still cannot compete with that of female plants, as their pod-bearing properties are worse. Future breeding would be focused on processing-quality aspects, as well as on properties for better mechanization of harvest or better-yielding hermaphroditic plants. The use of modern breeding techniques is restricted due to low polymorphism for molecular markers.
Uses
Food
Carob products consumed by humans come from the dried, sometimes roasted, pod, which has two main parts: the pulp accounts for 90% and the seeds 10% by weight. Carob pulp is sold either as flour or "chunks". The flour of the carob embryo (seed) can also be used for human and animal nutrition, but the seed is often separated before making carob powder (see section on locust bean gum below).
Carob pods are mildly sweet on their own (being roughly one third to one half sugar by dry weight), so they are used in powdered, chip or syrup form as an ingredient in cakes and cookies, sometimes as a substitute for chocolate in recipes because of the color, texture, and taste of carob. In Malta, a traditional sweet called karamelli tal-harrub and eaten during the Christian holidays of Lent and Good Friday is made from carob pods. Dried carob fruit is traditionally eaten on the Jewish holiday of Tu Bishvat.
Carob powder
Carob powder (carob pulp flour) is made of roasted, then finely ground, carob pod pulp.
Locust bean gum
Locust bean gum is produced from the endosperm, which accounts for 42–46% of the carob seed, and is rich in galactomannans (88% of endosperm dry mass). Galactomannans are hydrophilic and swell in water. If galactomannans are mixed with other gelling substances, such as carrageenan, they can be used to effectively thicken the liquid part of food. This is used extensively in canned food for animals in order to get the "jellied" texture.
Animal feed
While chocolate contains the chemical compound theobromine in levels that are toxic to some mammals, carob contains none, and it also has no caffeine, so it is sometimes used to make chocolate-like treats for dogs. Carob pod meal is also used as an energy-rich feed for livestock, particularly for ruminants, though its high tannin content may limit this use.
Historically, carob pods were mainly used for animal fodder in the Maltese islands, apart from times of famine or war, when they formed part of the diet of many Maltese people. On the Iberian Peninsula, carob pods were historically fed to donkeys.
Composition
The pulp of a carob pod is about 48–56% sugars and 18% cellulose and hemicellulose. Some differences in sugar (sucrose) content are seen between wild and cultivated carob trees: ~531 g/kg dry weight in cultivated varieties and ~437 g/kg in wild varieties. Fructose and glucose levels do not differ between cultivated and wild carob. The embryo (20-25% of seed weight) is rich in proteins (50%). The testa, or seed coat (30–33% of seed weight), contains cellulose, lignins, and tannins.
Syrup and drinks
Carob pods are about a third to a half sugar by weight, and this sugar can be extracted into a syrup. In Malta, a carob syrup (ġulepp tal-ħarrub) is made out of the pods. Carob syrup is also used in Crete, and Cyprus exports it.
In Egypt and Palestine, crushed pods are heated to caramelize their sugar, then water is added and boiled for some time. The result is a cold beverage, also called kharrub, which is sold by juice shops and street vendors, especially in summer. This drink is popular during Ramadan in Gaza, and in summer in Egypt.
In Lebanon the molasses is called debs el kharrub (literally: molasses of the carob), but people generally shorten it to debs. The molasses has a sweet, chocolate-like flavor. It is commonly mixed with tahini (typically 75% kharrub molasses and 25% tahini). The resulting mixture is called debs bi tahini and is eaten raw or with bread. The molasses is also used in certain cakes. The region of Iqlim al-Kharrub, which translates to the region of the carob, produces a significant amount of carob.
In Cyprus, the dried and milled carob pods are left to soak in water, before being transferred into special containers out of which the carob juice gradually seeps out of and is collected. The juice is then boiled with constant stirring yielding a thick syrup known as haroupomelo. Although this syrup is frequently sold and eaten as is, haroupomelo is also used as a base for a local toffee-like sweet snack known as pasteli. Constant stirring of the carob syrup causes it to form into a black, amorphous mass which is then left to cool. The mass is then kneaded, stretched and pulled until the fair, golden color and toffee-like texture of pasteli is obtained.
Carob is used for compote, liqueur, and syrup in Turkey, Malta, Portugal, Spain, and Sicily. In Libya, carob syrup (called rub) is used as a complement to asida (made from wheat flour). The so-called "carob syrup" made in Peru is actually from the fruit of the Prosopis nigra tree. Because of its strong taste, carob syrup is sometimes flavored with orange or chocolate. In Yemen, carob tree is playing a role in controlling diabetes mellitus according to Yemeni folk medicine, and diabetics consume carob pods as a juice to lower their blood sugar levels.
Ornamental
The carob tree is widely cultivated in the horticultural nursery industry as an ornamental plant for Mediterranean climates and other temperate regions around the world, being especially popular in California and Hawaii. The plant develops a sculpted trunk and the form of an ornamental tree after being "limbed up" as it matures, otherwise it is used as a dense and large screening hedge. The plant is very drought tolerant as long as one does not care about the size of the fruit harvest, so can be used in xeriscape landscape design for gardens, parks, and public municipal and commercial landscapes.
Timber
In some areas of Greece, viz. Crete, carob wood is often used as a firewood. As it makes such excellent fuel, it is sometimes even preferred over oak or olive wood.
Because the much fluted stem usually shows heart rot, carob wood is rarely used for construction timber. However, it is sometimes sought for ornamental work--particularly for furniture design, as the natural shape of the trunk is well-suited to the task. Additionally, the extremely wavy grain of the wood gives carob wood exceptional resistance to splitting; thus, sections of Carob bole are suitable for chopping blocks for splitting wood.
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| Biology and health sciences | Other culinary fruits | Plants |
347123 | https://en.wikipedia.org/wiki/Lymph | Lymph | Lymph () is the fluid that flows through the lymphatic system, a system composed of lymph vessels (channels) and intervening lymph nodes whose function, like the venous system, is to return fluid from the tissues to be recirculated. At the origin of the fluid-return process, interstitial fluid—the fluid between the cells in all body tissues—enters the lymph capillaries. This lymphatic fluid is then transported via progressively larger lymphatic vessels through lymph nodes, where substances are removed by tissue lymphocytes and circulating lymphocytes are added to the fluid, before emptying ultimately into the right or the left subclavian vein, where it mixes with central venous blood.
Because it is derived from interstitial fluid, with which blood and surrounding cells continually exchange substances, lymph undergoes continual change in composition. It is generally similar to blood plasma, which is the fluid component of blood. Lymph returns proteins and excess interstitial fluid to the bloodstream. Lymph also transports fats from the digestive system (beginning in the lacteals) to the blood via chylomicrons.
Bacteria may enter the lymph channels and be transported to lymph nodes, where the bacteria are destroyed. Metastatic cancer cells can also be transported via lymph.
Etymology
The word lymph is derived from the name of the ancient Roman deity of fresh water, Lympha.
Structure
Lymph has a composition similar but not identical to that of blood plasma. Lymph that leaves a lymph node is richer in lymphocytes than blood plasma is. The lymph formed in the human digestive system called chyle is rich in triglycerides (fat), and looks milky white because of its lipid content.
Development
Blood supplies nutrients and important metabolites to the cells of a tissue and collects back the waste products they produce, which requires exchange of respective constituents between the blood and tissue cells. This exchange is not direct, but instead occurs through an intermediary called interstitial fluid, which occupies the spaces between cells. As the blood and the surrounding cells continually add and remove substances from the interstitial fluid, its composition continually changes. Water and solutes can pass between the interstitial fluid and blood via diffusion across gaps in capillary walls called intercellular clefts; thus, the blood and interstitial fluid are in dynamic equilibrium with each other.
Interstitial fluid forms at the arterial (coming from the heart) end of capillaries because of the higher pressure of blood compared to veins, and most of it returns to its venous ends and venules; the rest (up to 10%) enters the lymph capillaries as lymph. (Prior to entry, this fluid is referred to as the lymph obligatory load, or LOL, as the lymphatic system is effectively "obliged" to return it to the cardiovascular network.) The lymph when formed is a watery clear liquid with the same composition as the interstitial fluid. However, as it flows through the lymph nodes it comes in contact with blood, and tends to accumulate more cells (particularly, lymphocytes) and proteins.
Functions
Components
Lymph returns proteins and excess interstitial fluid to the bloodstream. Lymph may pick up bacteria and transport them to lymph nodes, where the bacteria are destroyed. Metastatic cancer cells can also be transported via lymph. Lymph also transports fats from the digestive system (beginning in the lacteals) to the blood via chylomicrons.
Circulation
Tubular vessels transport lymph back to the blood, ultimately replacing the volume lost during the formation of the interstitial fluid. These channels are the lymphatic channels, or simply lymphatics.
Unlike the cardiovascular system, the lymphatic system is not closed. In some amphibian and reptilian species, the lymphatic system has central pumps, called lymph hearts, which typically exist in pairs, but humans and other mammals do not have a central lymph pump. Lymph transport is slow and sporadic. Despite low pressure, lymph movement occurs due to peristalsis (propulsion of the lymph due to alternate contraction and relaxation of smooth muscle tissue), valves, and compression during contraction of adjacent skeletal muscle and arterial pulsation.
Lymph that enters the lymph vessels from the interstitial spaces usually does not flow backwards along the vessels because of the presence of valves. If excessive hydrostatic pressure develops within the lymph vessels, though, some fluid can leak back into the interstitial spaces and contribute to formation of edema.
The flow of lymph in the thoracic duct in an average resting person usually approximates 100ml per hour. Accompanied by another ~25ml per hour in other lymph vessels, the total lymph flow in the body is about 4 to 5 litres per day. This can be elevated several fold while exercising. It is estimated that without lymphatic flow, the average resting person would die within 24 hours.
Clinical significance
Histopathological examination of the lymph system is used as a screening tool for immune system analysis in conjunction with pathological changes in other organ systems and clinical pathology to assess disease status. Although histological assessment of the lymph system does not directly measure immune function, it can be combined with identification of chemical biomarkers to determine underlying changes in the diseased immune system.
As a growth medium
In 1907 the zoologist Ross Granville Harrison demonstrated the growth of frog nerve cell processes in a medium of clotted lymph. It is made up of lymph nodes and vessels.
In 1913, E. Steinhardt, C. Israeli, and R. A. Lambert grew vaccinia virus in fragments of tissue culture from guinea pig cornea grown in lymph.
| Biology and health sciences | Circulatory system | Biology |
347603 | https://en.wikipedia.org/wiki/Carpenter%20bee | Carpenter bee | Carpenter bees are species in the genus Xylocopa of the subfamily Xylocopinae. The genus includes some 500 bees in 31 subgenera. The common name "carpenter bee" derives from their nesting behavior; nearly all species burrow into hard plant material such as dead wood or bamboo. The main exceptions are species in the subgenus Proxylocopa, which dig nesting tunnels in suitable soil.
Characteristics
Many species in this enormous genus are difficult to tell apart; most species are all black, or primarily black with some yellow or white pubescence. Some differ only in subtle morphological features, such as details of the male genitalia. Males of some species differ confusingly from the females, being covered in greenish-yellow fur. The confusion of species arises particularly in the common names; in India, for example, the common name for any all-black species of Xylocopa is bhanvra (or bhomora - ভোমোৰা - in Assamese), and reports and sightings of bhanvra or bhomora are commonly misattributed to a European species, Xylocopa violacea; however, this species is found only in the northern regions of Jammu and Kashmir and Punjab, and most reports of bhanvra, especially elsewhere in India, refer to any of roughly 15 other common black Xylocopa species in the region, such as X. nasalis, X. tenuiscapa, or X. tranquebarorum.
Non-professionals commonly confuse carpenter bees with bumblebees; the simplest rule of thumb for telling them apart is that most carpenter bees have a shiny abdomen, whereas bumblebee abdomens are completely covered with dense hair. Males of some species of carpenter bees have a white or yellow face, unlike bumblebees, while females lack the bare corbicula of bumblebees; the hind leg is entirely hairy.
The wing venation is characteristic; the marginal cell in the front wing is narrow and elongated, and its apex bends away from the costa. The front wing has small stigma. When closed, the bee's short mandibles conceal the labrum. The clypeus is flat. Males of many species have much larger eyes than the females, which relates to their mating behavior.
In the United States, two eastern species, Xylocopa virginica and X. micans, occur. Three more species are primarily western in distribution, X. sonorina, X. tabaniformis orpifex, and X. californica. X. virginica is by far the more widely distributed species.
Ecological significance
In several species, the females live alongside their own daughters or sisters, creating a small social group. They use wood bits to form partitions between the cells in the nest. A few species bore holes in wood dwellings. Since the tunnels are near the surface, structural damage is generally minor or superficial. However, carpenter bee nests are attractive to woodpeckers, which may do further damage by drilling into the wood to feed on the bees or larvae.
Carpenter bees have short mouthparts and are important pollinators on some open-faced or shallow flowers; for some they even are obligate pollinators, for example the maypop (Passiflora incarnata) and Orphium, which are not pollinated by any other insects. They also are important pollinators of flowers with various forms of lids, such as Salvia species and some members of the Fabaceae. However many carpenter bees "rob" nectar by slitting the sides of flowers with deep corollae. Xylocopa virginica is one example of a species with such nectar robbing behavior. With their short labia the bees cannot reach the nectar without piercing the long-tubed flowers; they miss contact with the anthers and perform no pollination. In some plants, this reduces fruit and seed production, while others have developed defence mechanisms against nectar robbing. When foraging for pollen from some species with tubular flowers however, the same species of carpenter bees still achieve pollination, if the anthers and stigmata are exposed together.
Many Old World carpenter bees have a special pouch-like structure on the inside of their first metasomal tergite called the acarinarium where certain mites (Dinogamasus species) reside as commensals. The exact nature of the relationship is not fully understood, though in other bees that carry mites, they are beneficial, feeding either on fungi in the nest, or on other harmful mites.
Behavior
As a subfamily, they nest in a wide range of host plants, but any one species may show definite adaptations or preferences for particular groups of plants. Carpenter bees are traditionally considered solitary bees, though some species have simple social nests in which mothers and daughters may cohabit. Examples of this type of social nesting can be seen in the species Xylocopa sulcatipes and Xylocopa nasalis. When females cohabit, a division of labor between them occurs sometimes. In this type of nesting, multiple females either share in the foraging and nest laying, or one female does all the foraging and nest laying, while the other females guard.
Solitary species differ from social species. Solitary bees tend to be gregarious and often several nests of solitary bees are near each other. In solitary nesting, the founding bee forages, builds cells, lays the eggs, and guards. Normally, only one generation of bees live in the nest. Xylocopa pubescens is one carpenter bee species that can have both social and solitary nests.
Carpenter bees make nests by tunneling into wood, bamboo, and similar hard plant material such as peduncles, usually dead. They vibrate their bodies as they rasp their mandibles against hard wood, each nest having a single entrance which may have many adjacent tunnels. The entrance is often a perfectly circular hole measuring about on the underside of a beam, bench, or tree limb. Carpenter bees do not eat wood. They discard the bits of wood, or reuse particles to build partitions between cells. The tunnel functions as a nursery for brood and storage for the pollen/nectar upon which the brood subsists. The provision masses of some species are among the most complex in shape of any group of bees; whereas most bees fill their brood cells with a soupy mass and others form simple spheroidal pollen masses, Xylocopa species form elongated and carefully sculpted masses that have several projections which keep the bulk of the mass from coming into contact with the cell walls, sometimes resembling an irregular caltrop. The eggs are very large relative to the size of the female, and are some of the largest eggs among all insects. Carpenter bees can be timber pests, and cause substantial damage to wood if infestations go undetected for several years.
Two very different mating systems appear to be common in carpenter bees, and often this can be determined simply by examining specimens of the males of any given species. Species in which the males have large eyes are characterized by a mating system where the males either search for females by patrolling, or by hovering and waiting for passing females, which they then pursue. In the other mating system, the males often have very small heads, but a large, hypertrophied glandular reservoir in the mesosoma releases pheromones into the airstream behind the male while it flies or hovers. The pheromone advertises the presence of the male to females.
Male bees often are seen hovering near nests and will approach nearby animals. However, males are harmless, since they do not have a stinger. Female carpenter bees are capable of stinging, but they are docile and rarely sting unless caught in the hand or otherwise directly provoked.
Natural predators
Woodpeckers eat carpenter bees, as do various species of birds, such as shrikes and bee-eaters as well as some mammals such as ratels. Other predators include large mantises and predatory flies, particularly large robber-flies of the family Asilidae. Woodpeckers are attracted to the noise of the bee larvae and drill holes along the tunnels to feed on them.
Apart from outright predators, parasitoidal species of bee flies (e.g. Xenox) lay eggs in the entrance to the bee's nest and the fly maggots live off the bee larvae.
Species
Xylocopa abbotti (Cockerell, 1909)
Xylocopa abbreviata Hurd & Moure, 1963
Xylocopa acutipennis Smith, 1854
Xylocopa adumbrata Lieftinck, 1957
Xylocopa adusta Pérez, 1901
Xylocopa aeneipennis (DeGeer, 1773)
Xylocopa aerata (Smith, 1851)
Xylocopa aestuans (Linnaeus, 1758)
Xylocopa aethiopica Pérez, 1901
Xylocopa africana (Fabricius, 1781)
Xylocopa albiceps Fabricius, 1804
Xylocopa albifrons Lepeletier, 1841
Xylocopa albinotum Matsumura, 1926
Xylocopa alternata Pérez, 1901
Xylocopa alticola (Cockerell, 1919)
Xylocopa amamensis Sonan, 1934
Xylocopa amauroptera Pérez, 1901
Xylocopa amazonica Enderlein, 1913
Xylocopa amedaei Lepeletier, 1841
Xylocopa amethystina (Fabricius, 1793)
Xylocopa andarabana Hedicke, 1938
Xylocopa andica Enderlein, 1913
Xylocopa angulosa Maa, 1954
Xylocopa anthophoroides Smith, 1874
Xylocopa apicalis Smith, 1854
Xylocopa appendiculata Smith, 1852
Xylocopa artifex Smith, 1874
Xylocopa aruana Ritsema, 1876
Xylocopa assimilis Ritsema, 1880
Xylocopa atamisquensis Lucia & Abrahamovich, 2010
Xylocopa augusti Lepeletier, 1841
Xylocopa auripennis Lepeletier, 1841
Xylocopa aurorea Friese, 1922
Xylocopa aurulenta (Fabricius, 1804)
Xylocopa bakeriana (Cockerell, 1914)
Xylocopa balteata Maa, 1943
Xylocopa bambusae Schrottky, 1902
Xylocopa bangkaensis Friese, 1903
Xylocopa barbatella Cockerell, 1931
Xylocopa bariwal Maidl, 1912
Xylocopa basalis Smith, 1854
Xylocopa bentoni Cockerell, 1919
Xylocopa bequaerti (Cockerell, 1930)
Xylocopa bhowara Maa, 1938
Xylocopa biangulata Vachal, 1899
Xylocopa bicarinata Alfken, 1932
Xylocopa bicristata Maa, 1954
Xylocopa bilineata Friese, 1914
Xylocopa bimaculata Friese, 1903
Xylocopa binongkona van der Vecht, 1953
Xylocopa bluethgeni Dusmet y Alonso, 1924
Xylocopa bombiformis Smith, 1874
Xylocopa bomboides Smith, 1879
Xylocopa bombylans (Fabricius, 1775)
Xylocopa boops Maidl, 1912
Xylocopa bouyssoui Vachal, 1898
Xylocopa brasilianorum (Linnaeus, 1767)
Xylocopa braunsi Dusmet y Alonso, 1924
Xylocopa bruesi Cockerell, 1914
Xylocopa bryorum (Fabricius, 1775)
Xylocopa buginesica Vecht, 1953
Xylocopa buruana Lieftinck, 1956
Xylocopa caerulea (Fabricius, 1804)
Xylocopa caffra (Linnaeus, 1767)
Xylocopa calcarata (LeVeque, 1928)
Xylocopa calens Lepeletier, 1841
Xylocopa californica Cresson, 1864
Xylocopa caloptera Pérez, 1901
Xylocopa canaria (Cockerell & LeVeque, 1925)
Xylocopa cantabrita Lepeletier, 1841
Xylocopa capensis Spinola, 1838
Xylocopa capitata Smith, 1854
Xylocopa carbonaria Smith, 1854
Xylocopa caribea Lepeletier, 1841
Xylocopa caspari van der Vecht, 1953
Xylocopa caviventris Maidl, 1912
Xylocopa cearensis Ducke, 1911
Xylocopa ceballosi Dusmet y Alonso, 1924
Xylocopa celebensis (Gribodo, 1894)
Xylocopa chapini (LeVeque, 1928)
Xylocopa chinensis Friese, 1911
Xylocopa chiyakensis (Cockerell, 1908)
Xylocopa chlorina (Cockerell, 1915)
Xylocopa chrysopoda Schrottky, 1902
Xylocopa chrysoptera Latreille, 1809
Xylocopa ciliata Burmeister, 1876
Xylocopa citrina Friese, 1909
Xylocopa clarionensis Hurd, 1958
Xylocopa claripennis Friese, 1922
Xylocopa cloti Vachal, 1898
Xylocopa cockerelli Maa, 1943
Xylocopa codinai Dusmet y Alonso, 1924
Xylocopa colona Lepeletier, 1841
Xylocopa columbiensis Pérez, 1901
Xylocopa combinata Ritsema, 1876
Xylocopa combusta Smith, 1854
Xylocopa concolorata Maa, 1938
Xylocopa conradsiana Friese, 1911
Xylocopa coracina van der Vecht, 1953
Xylocopa cornigera Friese, 1909
Xylocopa coronata Smith, 1861
Xylocopa cribrata Pérez, 1901
Xylocopa cubaecola Lucas, 1857
Xylocopa cuernosensis (Cockerell, 1915)
Xylocopa cyanea Smith, 1874
Xylocopa cyanescens Brullé, 1832
Xylocopa dalbertisi Lieftinck, 1957
Xylocopa dapitanensis (Cockerell, 1915)
Xylocopa darwini Cockerell, 1926
Xylocopa dejeanii Lepeletier, 1841
Xylocopa dibongoana Hedicke, 1923
Xylocopa dimidiata Latreille, 1809
Xylocopa disconota Friese, 1914
Xylocopa distinguenda Pérez, 1901
Xylocopa ditypa Vachal, 1898
Xylocopa diversipes Smith, 1861
Xylocopa dolosa Vachal, 1899
Xylocopa dormeyeri (Enderlein, 1909)
Xylocopa duala Strand, 1921
Xylocopa electa Smith, 1874
Xylocopa elegans Hurd & Moure, 1963
Xylocopa erlangeri Enderlein, 1903
Xylocopa erythrina Gribodo, 1894
Xylocopa escalerai Dusmet y Alonso, 1924
Xylocopa esica Cameron, 1902
Xylocopa euchlora Pérez, 1901
Xylocopa euxantha Cockerell, 1933
Xylocopa eximia Pérez, 1901
Xylocopa fabriciana Moure, 1960
Xylocopa fallax Maidl, 1912
Xylocopa fenestrata (Fabricius, 1798)
Xylocopa fervens Lepeletier, 1841
Xylocopa fimbriata Fabricius, 1804
Xylocopa flavicollis (DeGeer, 1778)
Xylocopa flavifrons Matsumura, 1912
Xylocopa flavonigrescens Smith, 1854
Xylocopa flavorufa (DeGeer, 1778)
Xylocopa forbesii W. F. Kirby, 1883
Xylocopa forsiusi Dusmet y Alonso, 1924
Xylocopa fortissima Cockerell, 1930
Xylocopa fransseni van der Vecht, 1953
Xylocopa friesiana Maa, 1939
Xylocopa frontalis (Olivier, 1789)
Xylocopa fuliginata Pérez, 1901
Xylocopa fulva Friese, 1922
Xylocopa funesta Maidl, 1912
Xylocopa fuscata Smith, 1854
Xylocopa gabonica (Gribodo, 1894)
Xylocopa gabrielae Engel, 2001
Xylocopa ganglbaueri Maidl, 1912
Xylocopa gaullei Vachal, 1898
Xylocopa ghilianii Gribodo, 1891
Xylocopa gracilis Dusmet y Alonso, 1923
Xylocopa graueri Maidl, 1912
Xylocopa gressitti Lieftinck, 1957
Xylocopa gribodoi Magretti, 1892
Xylocopa grisescens Lepeletier, 1841
Xylocopa griswoldi Mérida, Hinojosa-Díaz, & Ayala, 2022
Xylocopa grossa (Drury, 1770)
Xylocopa grubaueri Friese, 1903
Xylocopa gualanensis Cockerell, 1912
Xylocopa guatemalensis Cockerell, 1912
Xylocopa guigliae Lieftinck, 1957
Xylocopa haefligeri Friese, 1909
Xylocopa haematospila Moure, 1951
Xylocopa hafizii Maa, 1938
Xylocopa hellenica Spinola, 1843
Xylocopa hirsutissima Maidl, 1912
Xylocopa hottentotta Smith, 1854
Xylocopa hyalinipennis Friese, 1922
Xylocopa ignescens (LeVeque, 1928)
Xylocopa imitator Smith, 1854
Xylocopa incandescens (Cockerell, 1932)
Xylocopa incerta Pérez, 1901
Xylocopa incompleta Ritsema, 1880
Xylocopa inconspicua Maa, 1937
Xylocopa inconstans Smith, 1874
Xylocopa inquirenda Vachal, 1899
Xylocopa insola Vachal, 1910
Xylocopa insularis Smith, 1857
Xylocopa io Vachal, 1898
Xylocopa iranica Maa, 1954
Xylocopa iridipennis Lepeletier, 1841
Xylocopa iris (Christ, 1791)
Xylocopa isabelleae Hurd, 1959
Xylocopa javana Friese, 1914
Xylocopa kamerunensis Vachal, 1899
Xylocopa karnyi Maidl, 1912
Xylocopa kerri (Cockerell, 1929)
Xylocopa kuehni Friese, 1903
Xylocopa lachnea Moure, 1951
Xylocopa lanata Smith, 1854
Xylocopa langi (LeVeque, 1928)
Xylocopa lateralis Say, 1837
Xylocopa lateritia Smith, 1854
Xylocopa laticeps
Xylocopa latipes (Drury, 1773)
Xylocopa lautipennis (Cockerell, 1933)
Xylocopa lehmanni Friese, 1903
Xylocopa lepeletieri Enderlein, 1903
Xylocopa leucocephala Ritsema, 1876
Xylocopa leucothoracoides Maidl, 1912
Xylocopa levequeae Maa, 1943
Xylocopa lieftincki Leys, 2000
Xylocopa lombokensis Maidl, 1912
Xylocopa longespinosa Enderlein, 1903
Xylocopa longula Friese, 1922
Xylocopa loripes Smith, 1874
Xylocopa lucbanensis (Cockerell, 1927)
Xylocopa lucida Smith, 1874
Xylocopa lugubris Gerstäcker, 1857
Xylocopa lundqvisti Lieftinck, 1957
Xylocopa luteola Lepeletier, 1841
Xylocopa macrops Lepeletier, 1841
Xylocopa madida Friese, 1925
Xylocopa madurensis Friese, 1913
Xylocopa maesoi Dusmet y Alonso, 1924
Xylocopa magnifica (Cockerell, 1929)
Xylocopa maidli Maa, 1940
Xylocopa maior Maidl, 1912
Xylocopa marginella Lepeletier, 1841
Xylocopa mastrucata Pérez, 1901
Xylocopa maya Mérida, Hinojosa-Díaz, & Ayala, 2022
Xylocopa mazarredoi Dusmet y Alonso, 1924
Xylocopa mcgregori Cockerell, 1920
Xylocopa mckeani (Cockerell, 1929)
Xylocopa meadewaldoi Hurd, 1959
Xylocopa mendozana Enderlein, 1913
Xylocopa merceti Dusmet y Alonso, 1924
Xylocopa metallica Smith, 1874
Xylocopa mexicanorum Cockerell, 1912
Xylocopa meyeri Dusmet y Alonso, 1924
Xylocopa micans Lepeletier, 1841
Xylocopa micheneri Hurd, 1978
Xylocopa mimetica Cockerell, 1915
Xylocopa minor Maidl, 1912
Xylocopa mirabilis Hurd & Moure, 1963
Xylocopa mixta Radoszkowski, 1881
Xylocopa modesta Smith, 1854
Xylocopa mohnikei Cockerell, 1907
Xylocopa mongolicus (Wu, 1983)
Xylocopa montana Enderlein, 1903
Xylocopa mordax Smith, 1874
Xylocopa morotaiana Lieftinck, 1956
Xylocopa muscaria (Fabricius, 1775)
Xylocopa myops Ritsema, 1876
Xylocopa nasalis Westwood, 1842
Xylocopa nasica Pérez, 1901
Xylocopa nautlana Cockerell, 1904
Xylocopa negligenda Maa, 1939
Xylocopa nigrella Hurd, 1959
Xylocopa nigrescens Friese, 1901
Xylocopa nigricans Vachal, 1910
Xylocopa nigricaula (LeVeque, 1928)
Xylocopa nigripes Friese, 1915
Xylocopa nigrita (Fabricius, 1775)
Xylocopa nigrocaerulea Smith, 1874
Xylocopa nigrocaudata Pérez, 1901
Xylocopa nigrocincta Smith, 1854
Xylocopa nigroclypeata Rayment, 1935
Xylocopa nigroplagiata Ritsema, 1876
Xylocopa nigrotarsata Maa, 1938
Xylocopa nitidiventris Smith, 1878
Xylocopa nix (Maa, 1954)
Xylocopa nobilis Smith, 1859
Xylocopa nogueirai Hurd & Moure, 1960
Xylocopa nyassica Enderlein, 1903
Xylocopa oblonga Smith, 1874
Xylocopa obscurata Smith, 1854
Xylocopa obscuritarsis Friese, 1922
Xylocopa occipitalis Pérez, 1901
Xylocopa ocellaris Pérez, 1901
Xylocopa ocularis Pérez, 1901
Xylocopa ogasawarensis Matsumura, 1932
Xylocopa olivacea (Fabricius, 1778)
Xylocopa olivieri Lepeletier, 1841
Xylocopa ordinaria Smith, 1874
Xylocopa ornata Smith, 1874
Xylocopa orthogonaspis Moure, 2003
Xylocopa orthosiphonis (Cockerell, 1908)
Xylocopa pallidiscopa Hurd, 1961
Xylocopa parviceps Morawitz, 1895
Xylocopa parvula Rayment, 1935
Xylocopa perforator Smith, 1861
Xylocopa perkinsi Cameron, 1901
Xylocopa perpunctata (LeVeque, 1928)
Xylocopa peruana Pérez, 1901
Xylocopa perversa Wiedemann, 1824
Xylocopa pervirescens Cockerell, 1931
Xylocopa phalothorax Lepeletier, 1841
Xylocopa philippinensis Smith, 1854
Xylocopa pilosa Friese, 1922
Xylocopa plagioxantha Lieftinck, 1964
Xylocopa praeusta Smith, 1854
Xylocopa prashadi Maa, 1938
Xylocopa preussi Enderlein, 1903
Xylocopa provida Smith, 1863
Xylocopa proximata Maa, 1938
Xylocopa przewalskyi Morawitz, 1886
Xylocopa pseudoleucothorax Maidl, 1912
Xylocopa pseudoviolacea Popov, 1947
Xylocopa pubescens Spinola, 1838
Xylocopa pulchra Smith, 1874
Xylocopa punctifrons Cockerell, 1917
Xylocopa punctigena Maa, 1938
Xylocopa punctilabris Morawitz, 1894
Xylocopa pusulata Vachal, 1910
Xylocopa ramakrishnai Maa, 1938
Xylocopa rejecta Vachal, 1910
Xylocopa remota Maa, 1938
Xylocopa rogenhoferi Friese, 1900
Xylocopa romeroi Villamizar, Fernández, & Vivallo, 2020
Xylocopa rotundiceps Smith, 1874
Xylocopa rufa Friese, 1901
Xylocopa ruficeps Friese, 1910
Xylocopa ruficollis Hurd & Moure, 1963
Xylocopa ruficornis Fabricius, 1804
Xylocopa rufidorsum Enderlein, 1913
Xylocopa rufipes Smith, 1852
Xylocopa rufitarsis Lepeletier, 1841
Xylocopa rutilans Lieftinck, 1957
Xylocopa samarensis (Cockerell & LeVeque, 1925)
Xylocopa sarawatica Engel, 2017
Xylocopa schoana Enderlein, 1903
Xylocopa scioensis Gribodo, 1884
Xylocopa senex Friese, 1909
Xylocopa senior Vachal, 1899
Xylocopa shelfordi Cameron, 1902
Xylocopa sicheli Vachal, 1898
Xylocopa signata Morawitz, 1875
Xylocopa similis Smith, 1874
Xylocopa simillima Smith, 1854
Xylocopa sinensis (Wu, 1983)
Xylocopa sinensis Smith, 1854
Xylocopa smithii Ritsema, 1876
Xylocopa sogdiana Popov & Ponomareva, 1961
Xylocopa somalica Magretti, 1895
Xylocopa sonorina Smith, 1874
Xylocopa sphinx Vachal, 1899
Xylocopa splendidula Lepeletier, 1841
Xylocopa stadelmanni Vachal, 1899
Xylocopa stanleyi (LeVeque, 1928)
Xylocopa steindachneri Maidl, 1912
Xylocopa strandi Dusmet y Alonso, 1924
Xylocopa subcombusta (LeVeque, 1928)
Xylocopa subcyanea Pérez, 1901
Xylocopa subjuncta Vachal, 1898
Xylocopa subvirescens Cresson, 1879
Xylocopa subvolatilis (Cockerell, 1918)
Xylocopa subzonata Moure, 1949
Xylocopa sulcatipes Maa, 1970
Xylocopa sulcifrons Pérez, 1901
Xylocopa suspecta Moure & Camargo, 1988
Xylocopa suspiciosa Vachal, 1899
Xylocopa sycophanta Pérez, 1901
Xylocopa tabaniformis Smith, 1854
Xylocopa tacanensis Moure, 1949
Xylocopa tambelanensis (Cockerell, 1926)
Xylocopa tanganyikae Strand, 1911
Xylocopa tayabanica Cockerell, 1930
Xylocopa tegulata Friese, 1911
Xylocopa tenkeana Cockerell, 1933
Xylocopa tenuata Smith, 1874
Xylocopa tenuiscapa Westwood, 1840
Xylocopa teredo Guilding, 1825
Xylocopa tesselata Maa, 1970
Xylocopa thoracica Friese, 1903
Xylocopa togoensis Enderlein, 1903
Xylocopa torrida (Westwood, 1838)
Xylocopa tranquebarica (Fabricius, 1804)
Xylocopa tranquebarorum (Swederus, 1787)
Xylocopa transitoria Pérez, 1901
Xylocopa tricolor Ritsema, 1876
Xylocopa trifasciata Gribodo, 1891
Xylocopa trochanterica Vachal, 1910
Xylocopa truxali Hurd & Moure, 1963
Xylocopa tumida Friese, 1903
Xylocopa tumorifera Lieftinck, 1957
Xylocopa turanica Morawitz, 1875
Xylocopa uclesiensis Pérez, 1901
Xylocopa unicolor Smith, 1861
Xylocopa ustulata Smith, 1854
Xylocopa vachali Pérez, 1901
Xylocopa valga Gerstäcker, 1872
Xylocopa varentzowi Morawitz, 1895
Xylocopa varians Smith, 1874
Xylocopa varipes Smith, 1854
Xylocopa velutina Lieftinck, 1957
Xylocopa versicolor Alfken, 1930
Xylocopa vestita Hurd & Moure, 1963
Xylocopa villosa Friese, 1909
Xylocopa violacea (Linnaeus, 1758)
Xylocopa virginica (Linnaeus, 1771)
Xylocopa viridigastra Lepeletier, 1841
Xylocopa viridis Smith, 1854
Xylocopa vittata Enderlein, 1903
Xylocopa vogtiana Enderlein, 1913
Xylocopa volatilis Smith, 1861
Xylocopa vulpina Alfken, 1930
Xylocopa waterhousei Leys, 2000
Xylocopa watmoughi Eardley, 1983
Xylocopa wellmani Cockerell, 1906
Xylocopa wilmattae Cockerell, 1912
Xylocopa xanti Mocsáry, 1883
Xylocopa yunnanensis Wu, 1982
Xylocopa zonata Alfken, 1930
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| Biology and health sciences | Hymenoptera | Animals |
347838 | https://en.wikipedia.org/wiki/Nuclear%20medicine | Nuclear medicine | Nuclear medicine (nuclear radiology, nucleology), is a medical specialty involving the application of radioactive substances in the diagnosis and treatment of disease. Nuclear imaging is, in a sense, radiology done inside out, because it records radiation emitted from within the body rather than radiation that is transmitted through the body from external sources like X-ray generators. In addition, nuclear medicine scans differ from radiology, as the emphasis is not on imaging anatomy, but on the function. For such reason, it is called a physiological imaging modality. Single photon emission computed tomography (SPECT) and positron emission tomography (PET) scans are the two most common imaging modalities in nuclear medicine.
Diagnostic medical imaging
Diagnostic
In nuclear medicine imaging, radiopharmaceuticals are taken internally, for example, through inhalation, intravenously, or orally. Then, external detectors (gamma cameras) capture and form images from the radiation emitted by the radiopharmaceuticals. This process is unlike a diagnostic X-ray, where external radiation is passed through the body to form an image.
There are several techniques of diagnostic nuclear medicine.
2D: Scintigraphy ("scint") is the use of internal radionuclides to create two-dimensional images.
3D: SPECT is a 3D tomographic technique that uses gamma camera data from many projections and can be reconstructed in different planes. Positron emission tomography (PET) uses coincidence detection to image functional processes.
Nuclear medicine tests differ from most other imaging modalities in that nuclear medicine scans primarily show the physiological function of the system being investigated as opposed to traditional anatomical imaging such as CT or MRI. Nuclear medicine imaging studies are generally more organ-, tissue- or disease-specific (e.g.: lungs scan, heart scan, bone scan, brain scan, tumor, infection, Parkinson etc.) than those in conventional radiology imaging, which focus on a particular section of the body (e.g.: chest X-ray, abdomen/pelvis CT scan, head CT scan, etc.). In addition, there are nuclear medicine studies that allow imaging of the whole body based on certain cellular receptors or functions. Examples are whole body PET scans or PET/CT scans, gallium scans, indium white blood cell scans, MIBG and octreotide scans.
While the ability of nuclear metabolism to image disease processes from differences in metabolism is unsurpassed, it is not unique. Certain techniques such as fMRI image tissues (particularly cerebral tissues) by blood flow and thus show metabolism. Also, contrast-enhancement techniques in both CT and MRI show regions of tissue that are handling pharmaceuticals differently, due to an inflammatory process.
Diagnostic tests in nuclear medicine exploit the way that the body handles substances differently when there is disease or pathology present. The radionuclide introduced into the body is often chemically bound to a complex that acts characteristically within the body; this is commonly known as a tracer. In the presence of disease, a tracer will often be distributed around the body and/or processed differently. For example, the ligand methylene-diphosphonate (MDP) can be preferentially taken up by bone. By chemically attaching technetium-99m to MDP, radioactivity can be transported and attached to bone via the hydroxyapatite for imaging. Any increased physiological function, such as due to a fracture in the bone, will usually mean increased concentration of the tracer. This often results in the appearance of a "hot spot", which is a focal increase in radio accumulation or a general increase in radio accumulation throughout the physiological system. Some disease processes result in the exclusion of a tracer, resulting in the appearance of a "cold spot". Many tracer complexes have been developed to image or treat many different organs, glands, and physiological processes.
Hybrid scanning techniques
In some centers, the nuclear medicine scans can be superimposed, using software or hybrid cameras, on images from modalities such as CT or MRI to highlight the part of the body in which the radiopharmaceutical is concentrated. This practice is often referred to as image fusion or co-registration, for example SPECT/CT and PET/CT. The fusion imaging technique in nuclear medicine provides information about the anatomy and function, which would otherwise be unavailable or would require a more invasive procedure or surgery.
Practical concerns in nuclear imaging
Although the risks of low-level radiation exposures are not well understood, a cautious approach has been universally adopted that all human radiation exposures should be kept As Low As Reasonably Practicable, "ALARP". (Originally, this was known as "As Low As Reasonably Achievable" (ALARA), but this has changed in modern draftings of the legislation to add more emphasis on the "Reasonably" and less on the "Achievable".)
Working with the ALARP principle, before a patient is exposed for a nuclear medicine examination, the benefit of the examination must be identified. This needs to take into account the particular circumstances of the patient in question, where appropriate. For instance, if a patient is unlikely to be able to tolerate a sufficient amount of the procedure to achieve a diagnosis, then it would be inappropriate to proceed with injecting the patient with the radioactive tracer.
When the benefit does justify the procedure, then the radiation exposure (the amount of radiation given to the patient) should also be kept "ALARP". This means that the images produced in nuclear medicine should never be better than required for confident diagnosis. Giving larger radiation exposures can reduce the noise in an image and make it more photographically appealing, but if the clinical question can be answered without this level of detail, then this is inappropriate.
As a result, the radiation dose from nuclear medicine imaging varies greatly depending on the type of study. The effective radiation dose can be lower than or comparable to or can far exceed the general day-to-day environmental annual background radiation dose. Likewise, it can also be less than, in the range of, or higher than the radiation dose from an abdomen/pelvis CT scan.
Some nuclear medicine procedures require special patient preparation before the study to obtain the most accurate result. Pre-imaging preparations may include dietary preparation or the withholding of certain medications. Patients are encouraged to consult with the nuclear medicine department prior to a scan.
Analysis
The result of the nuclear medicine imaging process is a dataset comprising one or more images. In multi-image datasets the array of images may represent a time sequence (i.e. cine or movie) often called a "dynamic" dataset, a cardiac gated time sequence, or a spatial sequence where the gamma-camera is moved relative to the patient. SPECT (single photon emission computed tomography) is the process by which images acquired from a rotating gamma-camera are reconstructed to produce an image of a "slice" through the patient at a particular position. A collection of parallel slices form a slice-stack, a three-dimensional representation of the distribution of radionuclide in the patient.
The nuclear medicine computer may require millions of lines of source code to provide quantitative analysis packages for each of the specific imaging techniques available in nuclear medicine.
Time sequences can be further analysed using kinetic models such as multi-compartment models or a Patlak plot.
Interventional nuclear medicine
Radionuclide therapy can be used to treat conditions such as hyperthyroidism, thyroid cancer, skin cancer and blood disorders.
In nuclear medicine therapy, the radiation treatment dose is administered internally (e.g. intravenous or oral routes) or externally direct above the area to treat in form of a compound (e.g. in case of skin cancer).
The radiopharmaceuticals used in nuclear medicine therapy emit ionizing radiation that travels only a short distance, thereby minimizing unwanted side effects and damage to noninvolved organs or nearby structures. Most nuclear medicine therapies can be performed as outpatient procedures since there are few side effects from the treatment and the radiation exposure to the general public can be kept within a safe limit.
In some centers the nuclear medicine department may also use implanted capsules of isotopes (brachytherapy) to treat cancer.
History
The history of nuclear medicine contains contributions from scientists across different disciplines in physics, chemistry, engineering, and medicine. The multidisciplinary nature of nuclear medicine makes it difficult for medical historians to determine the birthdate of nuclear medicine. This can probably be best placed between the discovery of artificial radioactivity in 1934 and the production of radionuclides by Oak Ridge National Laboratory for medicine-related use, in 1946.
The origins of this medical idea date back as far as the mid-1920s in Freiburg, Germany, when George de Hevesy made experiments with radionuclides administered to rats, thus displaying metabolic pathways of these substances and establishing the tracer principle. Possibly, the genesis of this medical field took place in 1936, when John Lawrence, known as "the father of nuclear medicine", took a leave of absence from his faculty position at Yale Medical School, to visit his brother Ernest Lawrence at his new radiation laboratory (now known as the Lawrence Berkeley National Laboratory) in Berkeley, California. Later on, John Lawrence made the first application in patients of an artificial radionuclide when he used phosphorus-32 to treat leukemia.
Many historians consider the discovery of artificially produced radionuclides by Frédéric Joliot-Curie and Irène Joliot-Curie in 1934 as the most significant milestone in nuclear medicine. In February 1934, they reported the first artificial production of radioactive material in the journal Nature, after discovering radioactivity in aluminum foil that was irradiated with a polonium preparation. Their work built upon earlier discoveries by Wilhelm Konrad Roentgen for X-ray, Henri Becquerel for radioactive uranium salts, and Marie Curie (mother of Irène Curie) for radioactive thorium, polonium and coining the term "radioactivity." Taro Takemi studied the application of nuclear physics to medicine in the 1930s. The history of nuclear medicine will not be complete without mentioning these early pioneers.
Nuclear medicine gained public recognition as a potential specialty when on May 11, 1946, an article in the Journal of the American Medical Association (JAMA) by Massachusetts General Hospital's Dr. Saul Hertz and Massachusetts Institute of Technology's Dr. Arthur Roberts, described the successful use of treating Graves' Disease with radioactive iodine (RAI) was published. Additionally, Sam Seidlin. brought further development in the field describing a successful treatment of a patient with thyroid cancer metastases using radioiodine (I-131). These articles are considered by many historians as the most important articles ever published in nuclear medicine. Although the earliest use of I-131 was devoted to therapy of thyroid cancer, its use was later expanded to include imaging of the thyroid gland, quantification of the thyroid function, and therapy for hyperthyroidism. Among the many radionuclides that were discovered for medical-use, none were as important as the discovery and development of Technetium-99m. It was first discovered in 1937 by C. Perrier and E. Segre as an artificial element to fill space number 43 in the Periodic Table. The development of a generator system to produce Technetium-99m in the 1960s became a practical method for medical use. Today, Technetium-99m is the most utilized element in nuclear medicine and is employed in a wide variety of nuclear medicine imaging studies.
Widespread clinical use of nuclear medicine began in the early 1950s, as knowledge expanded about radionuclides, detection of radioactivity, and using certain radionuclides to trace biochemical processes. Pioneering works by Benedict Cassen in developing the first rectilinear scanner and Hal O. Anger's scintillation camera (Anger camera) broadened the young discipline of nuclear medicine into a full-fledged medical imaging specialty.
By the early 1960s, in southern Scandinavia, Niels A. Lassen, David H. Ingvar, and Erik Skinhøj developed techniques that provided the first blood flow maps of the brain, which initially involved xenon-133 inhalation; an intra-arterial equivalent was developed soon after, enabling measurement of the local distribution of cerebral activity for patients with neuropsychiatric disorders such as schizophrenia. Later versions would have 254 scintillators so a two-dimensional image could be produced on a color monitor. It allowed them to construct images reflecting brain activation from speaking, reading, visual or auditory perception and voluntary movement. The technique was also used to investigate, e.g., imagined sequential movements, mental calculation and mental spatial navigation.
By the 1970s most organs of the body could be visualized using nuclear medicine procedures. In 1971, American Medical Association officially recognized nuclear medicine as a medical specialty. In 1972, the American Board of Nuclear Medicine was established, and in 1974, the American Osteopathic Board of Nuclear Medicine was established, cementing nuclear medicine as a stand-alone medical specialty.
In the 1980s, radiopharmaceuticals were designed for use in diagnosis of heart disease. The development of single photon emission computed tomography (SPECT), around the same time, led to three-dimensional reconstruction of the heart and establishment of the field of nuclear cardiology.
More recent developments in nuclear medicine include the invention of the first positron emission tomography scanner (PET). The concept of emission and transmission tomography, later developed into single photon emission computed tomography (SPECT), was introduced by David E. Kuhl and Roy Edwards in the late 1950s. Their work led to the design and construction of several tomographic instruments at the University of Pennsylvania. Tomographic imaging techniques were further developed at the Washington University School of Medicine. These innovations led to fusion imaging with SPECT and CT by Bruce Hasegawa from University of California, San Francisco (UCSF), and the first PET/CT prototype by D. W. Townsend from University of Pittsburgh in 1998.
PET and PET/CT imaging experienced slower growth in its early years owing to the cost of the modality and the requirement for an on-site or nearby cyclotron. However, an administrative decision to approve medical reimbursement of limited PET and PET/CT applications in oncology has led to phenomenal growth and widespread acceptance over the last few years, which also was facilitated by establishing 18F-labelled tracers for standard procedures, allowing work at non-cyclotron-equipped sites. PET/CT imaging is now an integral part of oncology for diagnosis, staging and treatment monitoring. A fully integrated MRI/PET scanner is on the market from early 2011.
Sources of radionuclides
99mTc is normally supplied to hospitals through a radionuclide generator containing the parent radionuclide molybdenum-99. 99Mo is typically obtained as a fission product of 235U in nuclear reactors, however global supply shortages have led to the exploration of other methods of production. About a third of the world's supply, and most of Europe's supply, of medical isotopes is produced at the Petten nuclear reactor in the Netherlands. Another third of the world's supply, and most of North America's supply, was produced at the Chalk River Laboratories in Chalk River, Ontario, Canada until its permanent shutdown in 2018.
The most commonly used radioisotope in PET, 18F, is not produced in a nuclear reactor, but rather in a circular accelerator called a cyclotron. The cyclotron is used to accelerate protons to bombard the stable heavy isotope of oxygen 18O. The 18O constitutes about 0.20% of ordinary oxygen (mostly oxygen-16), from which it is extracted. The 18F is then typically used to make FDG.
A typical nuclear medicine study involves administration of a radionuclide into the body by intravenous injection in liquid or aggregate form, ingestion while combined with food, inhalation as a gas or aerosol, or rarely, injection of a radionuclide that has undergone micro-encapsulation. Some studies require the labeling of a patient's own blood cells with a radionuclide (leukocyte scintigraphy and red blood cell scintigraphy). Most diagnostic radionuclides emit gamma rays either directly from their decay or indirectly through electron–positron annihilation, while the cell-damaging properties of beta particles are used in therapeutic applications. Refined radionuclides for use in nuclear medicine are derived from fission or fusion processes in nuclear reactors, which produce radionuclides with longer half-lives, or cyclotrons, which produce radionuclides with shorter half-lives, or take advantage of natural decay processes in dedicated generators, i.e. molybdenum/technetium or strontium/rubidium.
The most commonly used intravenous radionuclides are technetium-99m, iodine-123, iodine-131, thallium-201, gallium-67, fluorine-18 fluorodeoxyglucose, and indium-111 labeled leukocytes. The most commonly used gaseous/aerosol radionuclides are xenon-133, krypton-81m, (aerosolised) technetium-99m.
Policies and procedures
Radiation dose
A patient undergoing a nuclear medicine procedure will receive a radiation dose. Under present international guidelines it is assumed that any radiation dose, however small, presents a risk. The radiation dose delivered to a patient in a nuclear medicine investigation, though unproven, is generally accepted to present a very small risk of inducing cancer. In this respect it is similar to the risk from X-ray investigations except that the dose is delivered internally rather than from an external source such as an X-ray machine, and dosage amounts are typically significantly higher than those of X-rays.
The radiation dose from a nuclear medicine investigation is expressed as an effective dose with units of sieverts (usually given in millisieverts, mSv). The effective dose resulting from an investigation is influenced by the amount of radioactivity administered in megabecquerels (MBq), the physical properties of the radiopharmaceutical used, its distribution in the body and its rate of clearance from the body.
Effective doses can range from 6 μSv (0.006 mSv) for a 3 MBq chromium-51 EDTA measurement of glomerular filtration rate to 11.2 mSv (11,200 μSv) for an 80 MBq thallium-201 myocardial imaging procedure. The common bone scan with 600 MBq of technetium-99m MDP has an effective dose of approximately 2.9 mSv (2,900 μSv).
Formerly, units of measurement were:
the curie (Ci), equal to 3.7 × 1010 Bq, and also equal to 1.0 grams of radium (Ra-226);
the rad (radiation absorbed dose), now replaced by the gray; and
the rem (Röntgen equivalent man), now replaced by the sievert.
The rad and rem are essentially equivalent for almost all nuclear medicine procedures, and only alpha radiation will produce a higher Rem or Sv value, due to its much higher Relative Biological Effectiveness (RBE). Alpha emitters are nowadays rarely used in nuclear medicine, but were used extensively before the advent of nuclear reactor and accelerator produced radionuclides. The concepts involved in radiation exposure to humans are covered by the field of Health Physics; the development and practice of safe and effective nuclear medicinal techniques is a key focus of Medical Physics.
Regulatory frameworks and guidelines
Different countries around the world maintain regulatory frameworks that are responsible for the management and use of radionuclides in different medical settings. For example, in the US, the Nuclear Regulatory Commission (NRC) and the Food and Drug Administration (FDA) have guidelines in place for hospitals to follow. With the NRC, if radioactive materials aren't involved, like X-rays for example, they are not regulated by the agency and instead are regulated by the individual states. International organizations, such as the International Atomic Energy Agency (IAEA), have regularly published different articles and guidelines for best practices in nuclear medicine as well as reporting on emerging technologies in nuclear medicine. Other factors that are considered in nuclear medicine include a patient's medical history as well as post-treatment management. Groups like International Commission on Radiological Protection have published information on how to manage the release of patients from a hospital with unsealed radionuclides.
| Biology and health sciences | Fields of medicine | Health |
14664934 | https://en.wikipedia.org/wiki/Glock | Glock | Glock (; stylized as GLOCK) is a brand of polymer-framed, short-recoil-operated, striker-fired, locked-breech semi-automatic pistols designed and produced by Austrian manufacturer Glock Ges.m.b.H.
The firearm entered Austrian military and police service in 1982 after becoming the top performer in reliability and safety tests.
Glock pistols have become the company's most profitable line of products, and have been supplied to national armed forces, security agencies, and police forces in at least 48 countries. Glocks are also popular firearms among civilians for recreational and competition shooting, home- and self-defense, both in concealed or open carry.
History
The company's founder and head engineer, Gaston Glock (1929–2023), had no experience with firearms design or manufacture at the time his first pistol, the Glock17, was being prototyped. Glock had extensive experience in advanced synthetic polymers, which was instrumental in the company's design of the first commercially successful line of pistols with a polymer frame. Glock introduced ferritic nitrocarburizing into the firearms industry as an anticorrosion surface treatment for metal gun parts.
Development
In 1980, the Austrian Armed Forces announced that it would seek tenders for a new, modern duty pistol to replace their World War II–era Walther P38 handguns. The Federal Ministry of Defence of Austria formulated a list of 17 criteria for the new generation service pistol, including requirements that it would be self loading; fire the NATO-standard 9×19mm Parabellum round; the magazines were not to require any means of assistance for loading; be secure against accidental discharge from shock, strike, and drop from a height of onto a steel plate. After firing 15,000 rounds of standard ammunition, the pistol was to be inspected for wear. The pistol was to then be used to fire an overpressure test cartridge generating . The normal maximum operating pressure (Pmax) for the 9 mm NATO is .
Glock became aware of the Austrian Army's planned procurement, and in 1982, assembled a team of Europe's leading handgun experts from military, police, and civilian sport-shooting circles to define the most desirable characteristics in a combat pistol. Within three months, Glock had developed a working prototype that combined proven mechanisms and traits from previous pistol designs. In addition, the plan was to make extensive use of synthetic materials and modern manufacturing technologies, which led to the Glock17 becoming a cost-effective candidate.
Several samples of the Glock17 (so named corresponding to its magazine capacity or because it was the 17th patent procured by the company) were submitted for assessment trials in early 1982, and after passing all of the exhaustive endurance and abuse tests, the Glock emerged as the winner. According to Friedrich Dechant, former head of the Austrian Armaments and Defence Technology Agency, the Glock P80 was clearly superior to other handguns in terms of performance, handling, charging capacity and price.
The handgun was adopted into service with the Austrian military and law enforcement in 1982 as the Pistole 80 (P80), with an initial order for 25,000 guns. The Glock17 outperformed eight different pistols from five other established manufacturers (Heckler & Koch of Germany offered their P7M8, P7M13, and P9S, SIG Sauer of Switzerland bid with their P220 and P226 models, Beretta of Italy submitted their model 92SB-F, FN Herstal of Belgium proposed an updated variant of the Browning Hi-Power, and the Austrian Steyr Mannlicher entered the competition with the GB).
The results of the Austrian trials sparked a wave of interest in Western Europe and overseas, particularly in the United States, where a similar effort to select a service-wide replacement for the M1911 had been going on since the late 1970s (known as the Joint Service Small Arms Program). In late 1983, the United States Department of Defense inquired about the Glock pistol and received four samples of the Glock17 for unofficial evaluation. Glock was then invited to participate in the XM9 Personal Defense Pistol Trials, but declined because the DOD specifications would require extensive retooling of production equipment and providing 35 test samples in an unrealistic time frame.
In 1985, after joint Norwegian and Swedish trials from 1983 to 1985, the Glock17 was accepted into service as the P80 in Norway, and in 1988 as the Pistol 88 in Sweden, where it surpassed all prior NATO durability standards. As a result, the Glock17 became a standard NATO-classified sidearm and was granted a NATO Stock Number (1005-25-133-6775). By 1992, some 350,000pistols had been sold in more than 45countries, including 250,000 in the United States alone.
Starting in 2013, the British Armed Forces began replacing the Browning Hi-Power pistol with the Glock17Gen4, due to concerns about weight and the external safety of the Hi-Power. The British preferred the Glock17Gen4 over the Beretta Px4 Storm, FN FNP, Heckler & Koch P30, SIG Sauer P226, Smith & Wesson M&P, and Steyr M9A1 of which 19 pistols each, all chambered in 9×19mmParabellum, were entered in the R9GSP trials.
The French Armed Forces (FAF) in 2020 began replacing their MAC Mle 1950 and, to a lesser extent, their PAMAS G1 pistols with Glock17Gen5 models specifically made for the FAF. The French preferred the Glock17Gen5 over the HS2000 and CZ P-10 offerings that also made it to the final selection phase.
Product evolution
Glock has updated its basic design several times throughout its production history.
First-generation models
The first-generation (Gen 1) Glock pistols are most notably recognized by their smoother "pebble finish" grip and finger groove-less frames. The Gen 1 frame pattern and design was used by Glock from 1982 through 1988 and pre-dates the checkered grip patterns used in the second generation of Glock pistols. The first Glock17s imported to the US were serialized with an alphanumeric (two-letter prefix followed by three numbers) stamped into the slide, barrel, and a small metal plate inserted into the bottom side of the polymer frame. The first documented Glock17s (by serial number) imported into the US were from the AF000series in January 1986, followed by AH000, AK000, and AL000. These early Glock(Gen 1) pistols (serial number prefixAF through AM) were also manufactured with a barrel that had a smaller overall diameter and thinner bore walls, later known as "pencil barrels". The barrels were later redesigned with thicker bore walls, and manufacturing continued to evolve and improve the design of Glock pistols.
Many of the first-generation Glocks were shipped and sold in the iconic "Tupperware" style plastic boxes. The earliest Glock boxes had ammunition storage compartments that allowed for 17 rounds of 9mm to be stored with the pistol. This box design was later changed by Glock to meet BATF import requirements, and the ammunition storage compartments were removed.
Second-generation models
A mid-life upgrade to the Glock pistols involved the addition of checkering on the front strap and trigger guard and checkering and serrations to the back strap. These versions, introduced in 1988, were informally referred to as "second-generation" or Gen2 models, though Glock did not mark the pistols Gen2. In 1991, an integrated recoil spring assembly replaced the original two-piece recoil spring and tube design. The magazine was slightly modified, changing the floorplate and fitting the follower spring with a resistance insert at its base.
Third-generation models
In 1998, the frame was further modified with an accessory rail (called the "Universal Glock rail") similar to a Picatinny rail to allow the mounting of laser sights, tactical lights, and other accessories. Thumb rests on both sides of the frame and finger grooves on the front strap were added. Glock pistols with these upgrades are informally referred to as (early) "third-generation" models. Later third-generation models additionally featured a modified extractor that serves as a loaded chamber indicator, and the locking block was enlarged, along with the addition of an extra cross pin to aid the distribution of bolt thrust forces exerted by the locking block. This cross pin is known as the locking block pin and is located above the trigger pin.
The polymer frames of third-generation models can be black, flat dark earth, or olive drab. Besides that, non-firing dummy pistols ("P"models) and non-firing dummy pistols with resetting triggers ("R"models) have a bright red frame, and Simunition-adapted practice pistols ("T"models) a bright blue frame for easy identification.
In 2009, the Glock22RTF2 (Rough Textured Frame 2) (chambered in .40 S&W) was introduced. This pistol featured a new checkering texture around the grip and new scalloped (fish gill-shaped) serrations at the rear of the sides of the slide. Many of the existing models became available in the RTF2version, including the 17, 31, 32, 23, 21, and 19. Some of those did not have the fish gills.
Fourth-generation models
At the 2010 SHOT Show, Glock presented the "fourth generation", now dubbed "Gen4" by Glock itself. Updates centered on ergonomics and the recoil spring assembly. The initial two fourth-generation models announced were the full-sized Glock 17 and Glock 22, chambered for the 9×19mm Parabellum and .40 S&W cartridges, respectively. The pistols were displayed with a modified rough-textured frame (RTF-4), front grip strap with finger grooves, interchangeable backstraps of different sizes, and an accessory rail. "Gen4" is rollmarked on the slide next to the model number to identify the fourth-generation pistols.
The basic grip size of the fourth-generation Glock pistols is slightly smaller compared to the previous design. A punch is provided to remove the standard trigger housing pin and replace it with the longer cross pin needed to mount the medium or large backstrap that will increase the trigger distance by or . With the medium backstrap installed, the grip size is identical to the third-generation pistols. The magazine release catches are enlarged and reversible for left-handed use. To use the exchangeable magazine release feature, fourth-generation Glock magazines have a notch cut on both sides of the magazine body. Earlier versions of the magazines will not lock into the Gen4 pistols if the user has moved the magazine release button to be operated by a left-handed user. Gen4 magazines will work in older models.
Mechanically, fourth-generation Glock pistols are fitted with a dual recoil spring assembly to help reduce perceived recoil and increase service life expectancy. Earlier subcompact Glock models such as the Glock 26 and Glock 30 have already used a dual recoil spring assembly that was carried over to the fourth-generation versions of those models. The slide and barrel shelf have been resized, and the front portion of the polymer frame has been widened and internally enlarged, to accommodate the dual recoil spring assembly. The trigger mechanism housing has also been modified to fit into the smaller-sized grip space.
The introduction of fourth-generation Glock pistols continued in July 2010 when the Glock 19 and Glock 23, the reduced size "compact" versions of the Glock 17 and Glock 22, became available for retail. In late 2010, Glock continued the introduction of fourth-generation models with the Glock 26 and Glock 27 "subcompact" variants.
In January 2013, more fourth-generation Glock pistols were introduced commercially during the annual SHOT Show, including the Glock 20 Generation 4 along with other fourth-generation Glock models.
2011 recoil spring assembly exchange program
In September 2011, Glock announced a recoil spring exchange program in which the manufacturer voluntarily offers to exchange the recoil spring assemblies of its fourth-generation pistols (with the exception of the "subcompact" Glock 26 and Glock 27 models) sold before 22 July 2011 at no cost "to ensure our products perform up to GLOCK's stringent standards", according to the company.
M series
On 29 June 2016, the United States Federal Bureau of Investigation (FBI) awarded a contract to Glock to provide new 9×19mm Parabellum chambered duty pistols. The solicitation specifications deviated from the specifications of Glock fourth-generation models. Features found in the M series pistols evolved into what become Glock's Fifth Generation or Gen5 pistols.
In August 2016, the Indianapolis Metro Police Department (IMPD) started training with a batch of Glock 17M pistols. The most obvious difference with the Glock third- and fourth-generation models on published images is the omission of finger grooves on the grip. In October of that year, the IMPD issued a 17M voluntary recall following failures encountered while dry firing the pistols during training. According to Major Riddle with the IMPD, "Glock is working to correct the problem and we hope to begin issuing the new [17Ms] as soon as December."
Fifth-generation models
In August 2017, Glock presented the "fifth generation" or "Gen 5". The revisions centered on ergonomics and improving reliability. Many parts of fifth-generation Glock pistols cannot be interchanged with those of the previous generations. The two fifth-generation models announced were the Glock 17 and Glock 19, chambered for the 9×19mm Parabellum. Some conspicuous changes on the fifth-generation models are ambidextrous slide stop levers, DLC surface finish for barrel and slide, a barrel featuring a revised style of polygonal rifling (called the "Glock Marksman Barrel" by Glock), a deeper recessed barrel crown, omission of the finger grooves on the grip, a flared magazine well, and a reintroduction of a half-moon-shaped cutout on the bottom front of the grip. The locking block pin located above the trigger pin that was introduced in the third generation is omitted. Many internal parts were less conspicuously revised. "Gen 5" is rollmarked on the slide next to the model number to identify the fifth-generation pistols. The "Gen 5" slide can feature front serrations (FS) to provide an additional tactile traction surface choice. The magazines were also revised for the fifth-generation models: the redesigned magazine floor plates feature a frontward protruding lip to offer grip for manual assisted extraction and the magazine follower became orange colored for easier visual identification.
Design details
Operating mechanism
The Glock 17 is a short recoil–operated, locked-breech semi-automatic pistol that uses a modified Browning cam-lock system adapted from the Hi-Power pistol. The firearm's locking mechanism uses a linkless, vertically tilting barrel with a rectangular breech that locks into the ejection port cut-out in the slide (the SIG Sauer system). During the recoil stroke, the barrel moves rearward initially locked together with the slide about until the bullet leaves the barrel and chamber pressure drops to a safe level. A ramped lug extension at the base of the barrel then interacts with a tapered locking block integrated into the frame, forcing the barrel down and unlocking it from the slide. This camming action terminates the barrel's movement while the slide continues back under recoil, extracting and ejecting the spent cartridge casing. The slide's uninterrupted rearward movement and counter-recoil cycle are characteristic of the Browning system.
Glock pistols incorporate a number of features intended to enhance reliability in adverse conditions, such as utilizing advanced metal coatings, "stub" slide guides instead of true frame rails, and an unusual cocking mechanism wherein the trigger is partially responsible for cocking the striker. By relying partially on force from the shooter's trigger finger to cock the striker, a Glock effectively reduces the load on the recoil spring as the slide moves forward into battery, whereas almost all other striker-fired pistols on the market rely fully on the recoil spring to cock the striker.
Features
The slide features a spring-loaded claw extractor, and the stamped sheet metal ejector is pinned to the trigger mechanism housing. Pistols after 2002 have a reshaped extractor that serves as a loaded chamber indicator. When a cartridge is present in the chamber, a tactile metal edge protrudes slightly out immediately behind the ejection port on the right side of the slide. The striker firing mechanism has a spring-loaded firing pin that is cocked in two stages that the firing pin spring powers. The factory-standard firing pin spring is rated at , but by using a modified firing pin spring, it can be increased to or to . When the pistol is charged, the firing pin is in the half-cock position. As the trigger is pulled, the firing pin is then fully cocked. At the end of its travel, the trigger bar is tilted downward by the connector, releasing the firing pin to fire the cartridge. The connector resets the trigger bar so that the firing pin will be captured in half-cock at the end of the firing cycle. This is known as a preset trigger mechanism, referred to as the "Safe action" trigger by the manufacturer. The connector ensures the pistol can only fire semiautomatically.
The factory-standard, two-stage trigger has a trigger travel of and is rated at , but by using a modified connector, it can be increased to or lowered to . In response to a request made by American law enforcement agencies for a two-stage trigger with increased trigger pull, Glock introduced the NY1 (New York) trigger module, which features a flat spring in a plastic housing that replaces the trigger bar's standard coil spring. This trigger modification is available in two versions: NY1 and NY2 that are rated at to and to , respectively, which require about to of force to disengage the safeties and another to in the second stage to fire a shot.
The Glock's frame, magazine body, and several other components are made from a high-strength nylon-based polymer invented by Gaston Glock, called Polymer 2. This plastic was specially formulated to provide increased durability and is more resilient than carbon steel and most steel alloys. Polymer 2 is resistant to shock, caustic liquids, and temperature extremes where traditional steel/alloy frames would warp and become brittle. The injection-molded frame contains four hardened steel guide rails for the slide: two at the rear of the frame, and the remaining pair above and in front of the trigger guard. The trigger guard itself is squared off at the front and checkered. The grip has an angle of 109° and a nonslip, stippled surface on the sides and both the front and rear straps. The frame houses the locking block, which is an investment casting that engages a 45° camming surface on the barrel's lower camming lug. It is retained in the frame by a steel axis pin that holds the trigger and slide catch. The trigger housing is held to the frame by means of a polymer pin. A spring-loaded sheet-metal pressing serves as the slide catch, which is secured from unintentional manipulation by a raised guard molded into the frame. Because of its polymer construction, there were initially fears that Glock pistols would be invisible to airport X-ray machines, making them easy to illegally import into the United States. In actuality, 84% of the gun's weight is from steel, and Polymer 2 is visible to X-ray machines. The myth's prevalence is believed to be connected to a scene that perpetuated the myth in Die Hard 2, which released a few years after the Glock was invented. In 1988, the Undetectable Firearms Act was passed in the United States, banning the manufacture or import of any gun that could pass undetected through a metal detector.
The Glock pistol has a relatively low slide profile, which holds the barrel axis close to the shooter's hand and makes the pistol more comfortable to fire by reducing muzzle rise and allows for faster aim recovery in rapid firing sequences. The rectangular slide is milled from a single block of ordnance-grade steel using CNC machinery. The barrel and slide undergo two hardening processes prior to treatment with a proprietary nitriding process called Tenifer. The Tenifer treatment is applied in a nitrate bath. The Tenifer finish is between in thickness, and is characterized by extreme resistance to wear and corrosion; it penetrates the metal, and treated parts have similar properties even below the surface to a certain depth.
The Tenifer process produces a matte gray-colored, nonglare surface with a 64 Rockwell C hardness rating and a 99% resistance to salt water corrosion (which meets or exceeds stainless steel specifications), making the Glock particularly suitable for individuals carrying the pistol concealed as the highly chloride-resistant finish allows the pistol to better endure the effects of perspiration. Glock steel parts using the Tenifer treatment are more corrosion resistant than analogous gun parts having other finishes or treatments, including Teflon, bluing, hard chrome plating, or phosphates. During 2010, Glock switched from the salt bath nitriding Tenifer process to a not exactly disclosed gas nitriding process. After applying the nitriding process, a black Parkerized decorative surface finish is applied. The underlying nitriding treatment will remain, protecting these parts even if the decorative surface finish were to wear off.
A fourth generation Glock 17 consists of 34 parts. For maintenance, the pistol disassembles into five main groups: the barrel, slide, frame, magazine, and recoil-spring assembly. The firearm is designed for the NATO-standard 9×19mm Parabellum cartridge, but can use high-power (increased pressure) +P ammunition with either full-metal-jacket or jacketed hollow-point projectiles.
Barrel
The hammer-forged barrel has a female type polygonal rifling with a right-hand twist. The stabilization of the round is not by conventional rifling, using lands and grooves, but rather through a polygonal profile consisting of a series of six or eight interconnected noncircular segments (only the .45 ACP and .45 GAP have octagonal polygonal rifling). Each depressed segment within the interior of the barrel is the equivalent of a groove in a conventional barrel. Thus, the interior of the barrel consists of smooth arcs of steel rather than sharply defined slots. Instead of using a traditional broaching machine to cut the rifling into the bore, the hammer forging process involves beating a slowly rotating mandrel through the bore to obtain the hexagonal or octagonal shape. As a result, the barrel's thickness in the area of each groove is not compromised as with conventional square-cut barrels. This has the advantage of providing a better gas seal behind the projectile as the bore has a slightly smaller diameter, which translates into more efficient use of the combustion gases trapped behind the bullet, slightly greater (consistency in) muzzle velocities, and increased accuracy and ease of maintenance.
The newer lines of Glock pistols—i.e. Gen5, G42/43—are equipped with the Glock Marksmanship Barrel, or GMB. While older barrels were somewhat difficult to identify a bullet as coming from a particular barrel with high enough reliability for evidentiary use, the newer GMB ones are designed differently. A study by Stephen Christen and Hans Rudolf Jordi, published by Forensic Science International in February 2019, shows that the new GMB barrels leave more identifiably unique markings on the fired projectile. These marks were more easily identified than previous pistol barrel markings, and were sufficient for reliably tying a bullet to a particular barrel. The study used a comparison microscope and an ABIS (Evofinder).
Safety
Glock pistols lack a traditional on-off safety lever, which Glock markets as an advantage, especially to police departments, as the user is able to fire immediately without separately manipulating a safety. Instead, the pistols are designed with three independent safety mechanisms to prevent accidental discharge. The system, designated "Safe Action" by Glock, consists of an external integrated trigger safety and two automatic internal safeties: a firing pin safety and a drop safety. The external safety is a small inner lever contained in the trigger. Pressing the lever activates the trigger bar and sheet metal connector. The firing pin safety is a solid hardened steel pin that, in the secured state, blocks the firing pin channel (disabling the firing pin in its longitudinal axis). It is pushed upward to release the firing pin for firing only when the trigger is actuated and the safety is pushed up through the backward movement of the trigger bar. The drop safety guides the trigger bar in a ramp that is released only when direct rearward pressure is applied to the trigger. The three safety mechanisms are automatically disengaged one after the other when the trigger is squeezed, and are automatically reactivated when the trigger is released.
In 2003, Glock announced the Internal Locking System (ILS) safety feature named Glock Safety Lock. The ILS is a manually activated lock located in the back of the pistol's grip. It is cylindrical in design and, according to Glock, each key is unique. When activated, the lock causes a tab to protrude from the rear of the grip, giving both a visual and tactile indication as to whether the lock is engaged or not. When activated, the ILS renders the Glock unfireable, as well as making it impossible to disassemble. When disengaged, the ILS adds no further safety mechanisms to the Glock pistol. The ILS is available as an option on most Glock pistols. Glock pistols cannot be retrofitted to accommodate the ILS. The lock must be factory-built in Austria and shipped as a special order.
Magazine
The Glock 17 feeds from staggered-column or double stack magazines that have a 17-round capacity (which can be extended to 19 with an optional floor plate) or optional 24 or 33-round high-capacity magazines. For jurisdictions which restrict magazine capacity to 10 rounds, Glock offers single-stack, 10-round magazines. The magazines are made of steel and are overmolded with plastic. A steel spring drives a plastic follower. After the last cartridge has been fired, the slide remains open on the slide stop. The slide stop release lever is located on the left side of the frame directly beneath the slide and can be manipulated by the thumb of the right-handed shooter.
Glock magazines are interchangeable between models of the same caliber, meaning that a compact or subcompact pistol will accept magazines designed for the larger pistols chambered for the same round. However, magazines designed for compact and subcompact models will not function in larger pistols because they are not tall enough to reach the slide and magazine release. For example, the subcompact Glock 26 will accept magazines from both the full-size Glock 17 and the compact Glock 19, but the Glock 17 will not accept magazines from the smaller Glock 19 or the Glock 26. The magazines for the Glock 36, the Glock 42, the Glock 43, and the Glock 44 are all unique; they cannot use magazines intended for another model, nor can their magazines be used in other models.
Sights
The first Glock pistols sent to the United States in 1985 failed to meet the BATF import "points" requirement, requiring Glock to quickly develop an adjustable rear sight which allowed for the pistols to be imported and sold commercially in 1986. It is believed that Glock designed and created this adjustable rear sight over a weekend in order to meet the ATF's importation requirements, and so it was dubbed the "weekend" sight. These first-generation adjustable rear sights extended past the slide and were susceptible to breaking. Even on later models, the front sight can easily become misshapen from friction against the holster, leading to replacements with metal sights, or tritium illuminated night sights.
More commonly today, the Glock 17 has a fixed polymer combat-type sighting arrangement that consists of a ramped front sight and a notched rear sight with white contrast elements painted on for increased acquisition speed – a white dot on the front post and a rectangular border on the rear notch. Some newer rear sights can be adjusted for windage (on certain models due to the windage sights not coming as factory default), as it has a degree of lateral movement in the dovetail it is mounted in. Three other factory rear sight configurations are available in addition to the standard height sight: a lower impact sight, and two higher impact versions – and .
Accessories
The Glock pistol accessories available from the factory include several devices for tactical illumination, such as a series of front rail-mounted "Glock tactical lights" featuring a white tactical light and an optional visible laser sight. An alternate version of the tactical light using an invisible infrared light and laser sight is available, designed to be used with an infrared night vision device. Another lighting accessory is an adapter to mount a flashlight onto the bottom of a magazine.
Polymer holsters in various configurations and matching magazine pouches are available. In addition, Glock produces optional triggers, recoil springs, slide stops, magazine release levers, and maritime spring cups. Maritime spring cups are designed to allow the pistol to be fired immediately after being submerged in water. They feature additional openings that allow liquids to flow and escape around them, offering enhanced reliability when water has penetrated into the firing pin assembly channel.
Magazine floor plates (or +2 baseplates), which expand the capacity of the standard magazines by two rounds, are available for models chambered for the 9×19mm Parabellum, .45 GAP, .40 S&W, .357 SIG, and .380 ACP cartridges. In addition to the standard nonadjustable polymer sight line, three alternative sight lines are offered by Glock. These consist of steel, adjustable, and self-illuminating tritium night rear sights and factory steel and self-illuminating tritium contrast pointer steel front sights.
The Glock 17 along with many variants can accept pistol conversion kits, with one such example being the FAB-Defense KPOS Scout. They can also accept special stocks like the Flux Defense Brace.
Glock switch
A Glock switch is an aftermarket accessory which depresses the firearm's sear, allowing fully automatic fire. Without the proper license, they are illegal in the United States.
Commemorative, anniversary, engraved, and other rare Glocks
Glock began producing limited edition and commemorative Glocks in 1991. Glock later produced a series of anniversary models to celebrate business milestones and in honor of 20, 25, and 30 years of US sales. Additionally, many law enforcement agencies had the department name, logo, or badges engraved on the slides of issued duty weapons.
Variants
Following the introduction of the Glock 17, numerous variants and versions have been offered. Variants that differ in caliber, frame, and slide length are identified by different model numbers with the exception of a few models with a letter suffix (the Glock 17L, 19X, 30S, and 43X).
The original double-stack "small frame" Glock pistols are made in five form factors, all modeled after the original full-sized Glock 17. "Standard" models are designed as full-sized duty firearms with a large magazine capacity. "Compact" models are slightly smaller with reduced magazine capacity and lighter weight while maintaining a usable grip length. "Subcompact" models are designed for easier carry and being lighter and shorter, are intended to be used with two fingers on the grip below the trigger guard and lack an accessory rail like the larger, after generation two, Glock models. The other two form factors use the full-size "standard" frame with longer slides that include a lower section to fill in the space between the frame's dust cover and the front of the slide. The first of these are the "long slide" models, which were too long for certain IPSC classes, necessitating the creation of the intermediate "competition" models. Currently, available chamberings for all five form-factors are 9mm Parabellum (9×21mm in certain countries) and .40 Smith & Wesson. Additionally, .357 SIG and .45 GAP chamberings are offered in "standard", "compact", and "subcompact" models, while .380 ACP is offered only in "compact" and "sub-compact" models. Recently, so-called "crossover" versions in 9mm Parabellum pair short (front to back) "compact" frames with longer "standard" grip lengths. This was initially to provide a longer-grip, higher capacity version of the Glock 19 (Glock 19X and Glock 45), but Glock developed the Glock 47 for US Customs and Border Patrol which used the G45 frame with a G17 length slide that included a front section to fill in the gap between the dust cover and the front of the slide, a design that fits the same overall dimensions as the Glock 17. Naturally, the next step was to couple this slide with the Glock 19 Frame, creating the Glock 49.
There are also the wider double-stack "large frame" Glock pistols for use with larger calibers, currently in 10mm Auto and .45 ACP. These models have bigger, wider slides and frames and are larger than the smaller-chambered pistols. These come in only "standard", "subcompact", "competition" (.45 ACP only), and "long slide" (10mm only). Additionally, Glock introduced the "subcompact" Glock 30S in .45 ACP adjusted to use the slimmer lighter slide of the Glock 36. Also, in 2007, Glock introduced a "short frame" version of these large frame weapons to provide a grip better suited to small hands. The short frame was originally designed to compete in the now cancelled U.S. military Joint Combat Pistol trials for a new .45 ACP pistol to replace the M9 pistol. Glock's entry featured an optional ambidextrous magazine release and MIL-STD-1913 rail along with a reduction in the size of the backstrap. The Glock 21SF was originally available in three versions: one with a Picatinny rail and ambidextrous magazine release and two with a Universal Glock rail available with or without the ambidextrous magazine release. However, the ambidextrous release and picatinny rail were soon dropped. As of January 2009, the Glock 20, 21, 29, and 30 were offered in short-framed variations. These models incorporate a reduction in trigger reach, and full-sized models feature a reduction in heel depth, which corresponds to an overall reduction in length for those models. The short frame models were not introduced for the Gen 4 and Gen 5 models, as the replaceable backstraps design makes a separate short frame version redundant.
Glock also produces single-stack "slimline" models, targeting the concealed carry market. The first was introduced with Gen 3 in .45 ACP as the Glock 36. More recently, after the introduction of Gen 5, came the Glock 42 in .380 ACP, followed by the Glock 43 in 9mm Parabellum. The most recent additions to the "slimline" series, the 43X and 48, were introduced together in 9mm Parabellum. These have longer grips that allow for a full three-finger hold and a 10-round capacity. The 43X is a long grip "crossover" 43, while the 48 has a longer slide to provide a "slimline" version of the "compact" Glock 19.
9×19mm Parabellum
Glock 17: The Glock 17 is the original 9×19mm Parabellum model, with a standard magazine capacity of 17 rounds, introduced in 1982. Initial samples of the new civilian offering were marked Glock 82. But, it was decided not to use the year to designate civilian models, but to begin the model numbers at 17 and continue numerically from there. Glock also offers a version of the standard magazine which incorporates a longer "+2" base plate to provide a capacity of 19 rounds. Also, a 10-round version of the standard magazine was created for markets that restrict the magazine capacity of handguns. Glock also offers an extended 24-round (with flush base plate) magazine for the Glock 17. Finally, the Glock 17 can use the Glock 18's extended 33-round (with +2 base plate) magazine. The base plates for the extended magazines can be swapped out to create 26- and 31-round magazines as well. The longest serving of the Glocks, the Glock 17 can be had with numerous "options" such as a threaded barrel, or slides cut for the Modular Optic System (MOS). Some options, such as the universal Glock rail have become standard. In addition, some features have been given their own suffixed model designations creating entirely new models, all of which can use the same magazines as the Glock 17:
Glock 17L: Introduced in 1988, the 17L incorporates a longer slide and extended barrel. Initially, the Glock 17L had three holes in the top of the barrel and a corresponding slot in the slide; however, later production pistols lack the holes in the barrel. The Glock 17L is manufactured in limited quantities.
Glock 17C: Introduced in 1996, the 17C incorporates slots cut in the barrel and slide to compensate for recoil. Many other Glock pistols now come with this option, all with a "C" suffix on the slide.
Glock 17MB: The 17MB is a version with ambidextrous magazine catch. This model, along with the other MB variants, was no longer available upon the introduction of the fourth-generation models, which have a reversible magazine catch.
Glock 17M: Introduced in 2016, the 17M was created in response to an FBI solicitation for a new full-size 9mm pistol. Differences from the Generation 4 model include removal of the finger grooves, ambidextrous slide lock, rounded slide nose profile, flared magazine well with new magazine baseplates, and a tougher finish on metal components. The Glock 17M also abandons the polygonal rifling of previous models for conventional rifling. , the Federal Bureau of Investigation, the South Carolina Highway Patrol and the Ontario Provincial Police have adopted the pistol as standard.
Glock P80: Introduced in 2020, the P80 was commissioned by United States firearms distributor Lipsey's to create an exclusive commemorative Glock model, the Pistole 80. The P80 is a throwback to the original Glock 17 Gen 1 type pistol chambered in 9×19mm with original Gen 1 frame and stippling and Gen 2 / Gen 3 internals.
Glock 18: The Glock 18 is a selective-fire variant of the Glock 17, developed at the request of the Austrian counter-terrorist unit EKO Cobra, and as a way to internally test Glock components under high strain conditions. Originally produced in 1986, this machine pistol–class firearm has a lever-type fire-control circular selector switch, installed on the serrated portion of the rear left side of the slide. With the selector lever in the bottom position, the pistol fires fully automatically at a cyclic rate of 1,100–1,400 RPM (rounds per minute), and with the selector lever in the top position, the pistol fires semi-automatically. The firearm is typically used with an extended 33-round-capacity magazine and may be fired with or without a shoulder stock, although other magazines from the Glock 17 can be used, with available capacities of 10, 17, 19 or 24 rounds. Unlike all its other pistols, it is only offered to military, law enforcement, and government organizations. Early Glock 18 models were ported to reduce muzzle rise during automatic fire. A very early design introduced a longer ported barrel which was soon discarded as it would not fit in a holster. Another compensated variant was produced, known as the Glock 18C. It has a keyhole opening cut into the forward portion of the slide, similar to the opening on the Glock long-slide models, although the Glock 18 has a standard-length slide. The keyhole opening provides an area to allow the four, progressively larger (from back to front) compensator cuts machined into the barrel to vent the propellant gases upwards, affording more control over the rapid-firing machine pistol.
Glock 18C: The compensator cuts start about halfway back on the top of the barrel. The two rear cuts are narrower than the two front cuts. The slide is hollowed, or dished-out, in a rectangular pattern between the rear of the ejection port and the rear sight. The rate of fire in fully automatic mode is around 1,100–1,200 rounds per minute. Most of the other characteristics are equivalent to the Glock 17, although the slide, frame, and certain fire-control parts of the Glock 18 are not interchangeable with other Glock models.
Glock 19: The Glock 19 is effectively a reduced-size Glock 17, called the "Compact" by the manufacturer. It was first produced in 1988, primarily for military and law enforcement. The Glock 19's barrel and pistol grip are shorter by about than the Glock 17, and it uses a magazine with a standard capacity of 15 rounds. A 10-round version of this magazine is also made for markets that restrict the magazine capacity of handguns. And, a "+2" base plate can make the standard magazine into a longer 17-round magazine. The pistol is also compatible with any magazines designed for the Glock 17 and Glock 18, providing factory magazine capacities of 17, 19, 24 and 33. Changing out base plates adds capacities of 26 or 31. To preserve the operational reliability of the short recoil system, the mass of the slide remains the same as in the Glock 17 from which it is derived. With the exception of the slide, frame, barrel, locking block, recoil spring, guide rod, and slide lock spring, all of the other components are interchangeable between the models 17 and 19. The Glock 19 Gen 4 MOS (Modular Optic System) has also been used by Special Operations Forces as the MK27 MOD 2. One of the oldest of the Glock pistols, options like threaded barrel and MOS slide cuts are available, and suffixed model designations have been created for some features, all of which excepting the G19X can use any magazine the G19 can:
Glock 19X: The 19X is the civilian version of Glock's entry to the XM17 Modular Handgun System competition for the United States Armed Forces. It features a Glock 19 slide with a Glock 17-like frame in coyote color instead of the regular Glock black color. The frame includes a lanyard loop and a front lip in the magazine which purpose is to make changing magazines with gloves on easier, but this means the new Gen 5 17-round magazines cannot be used in the Glock 19X because the front lip will block the extended magazine floor plates from locking into the 19X's magazine well. This can be remedied by switching to a Gen 4–style magazine floor plate or with a factory +2 extension. The G19X can use any factory G17 magazine for Gen 4 and prior. It is only the Gen 5 17- and 19-round magazines that it can't accept. The 19X comes standard with night sights and includes one 17-round magazine and two 19-round magazines, all in coyote color. The Glock 19X has proven to be one of Glock's best selling pistols, with over 100,000 sold within 6 months of the 19X first being released.
Glock 19M: Introduced in 2016, the 19M was created in response to an FBI solicitation for a new compact 9mm pistol. Differences from the Generation 4 model include removal of the finger grooves, ambidextrous slide stop, rounded slide nose profile, flared magazine well with new magazine baseplates, and a tougher finish on metal components. The Glock 19M also abandons the polygonal rifling of previous models for conventional rifling. The US Marine Corps fielded the Glock 19M, designated as the M007, to CID (Criminal Investigation Division) and Marine One personnel.
Glock 19 Canadian: The limit for Restricted Class firearms in Canada is a 105mm barrel, so, due to its 102mm barrel, the standard Glock 19 is too short to be legal. Starting in 2017, a market specific Glock 19 has been sold in Canada with a 106mm barrel and a distinctive laser-engraved hollow maple leaf on the right side of the slide.
Glock 26: The Glock 26 is a 9×19mm "subcompact" variant designed for concealed carry and was introduced in 1995, mainly for the civilian market. It has also been acquired by the US military and designated MK 26. It features a smaller frame compared to the Glock 19, with a pistol grip that supports only two fingers, a shorter barrel and slide, and a double-stack magazine with a standard capacity of 10 rounds. A factory magazine with a +2 extension gives a capacity of 12 rounds. In addition, the Glock 26 can use factory magazines from the Glock 17, Glock 18, and Glock 19, and one can swap out base plates to give it capacities of 15, 17, 19, 24, 26, 31 and 33 rounds. More than simply a "shortened" Glock 19, design of the subcompact Glock 26 required extensive rework of the frame, locking block, and spring assembly that features a dual recoil spring. MOS options are available, but so far the Glock universal rail is not offered. There is also:
Glock 26 for U.S. Customs and Border Protection (CBP): This Glock is a Generation 5 Glock for the CBP that incorporates a flared magazine well with an extended, longer grip than that of the usual Gen 5 Glock 26. Moreover, the magazine well is flat across and has no bump as the Glock 25 Gen 5 or the G19X. In addition, it offers a longer 11-round magazine.
Glock 34: The Glock 34 is a competition version of the Glock 17. It is similar to its predecessor, the Glock 17L, but with a slightly shorter slide and barrel, to meet the maximum size requirements for many sanctioned action pistol sporting events. It was developed and produced in 1998, and compared to the Glock 17, features a longer barrel and slide. It has an extended magazine release, extended slide stop lever, trigger pull, and an adjustable rear sight. The sides at the front of the slide are slanted instead of squared. Further, the top of the slide and parts of its inside are milled out, creating a conspicuous hole at the top designed to reduce front-end muzzle weight to better balance the pistol and reduce the overall weight of the slide. The Glock 34 can accept any magazine the Glock 17 can accept.
Glock 43: The Glock 43 is a "slimline" version of the subcompact Glock 26 that features an ultracompact slide and frame. The Glock 43 is the first Glock pistol to be manufactured with a single-stack 9×19mm Parabellum magazine, having a standard capacity of six rounds and being unique to the model. Unlike other subcompact Glock pistols, the Glock 43 cannot use factory magazines from its larger relatives due to its single-stack magazine design. It also does not allow the removal of the backplate grip as is possible on the 4th gen Glocks. The magazine is thinner than the Glock 43X, and the Glock 43 cannot accept magazines for the Glock 43X.
Glock 43X: The 43X is similar to the 43 except it has a longer and thicker grip for an increased magazine capacity of 10 rounds. The grip of the Glock 43X is comparable to the 48 and can be interchanged. Glock 43X magazines do not fit into the Glock 43, or vice versa. The 43X also features front slide serrations, a built-in extended beaver tail, a reversible magazine catch (similar to Gen 5 models), GMB rifling (again similar to Gen 5 models), and a two-tone finish (silver slide/black receiver). In the E.U. the 43X comes with a rail. At least three aftermarket sources manufacture 15-round flush-fit magazines for the Glock 43X and Glock 48, which make the Glock 43X and Glock 48 match the standard capacity of the Glock 19 in a narrower pistol.
Glock 45: The Glock 45 , similar to the Glock 19X, incorporates Gen5 features catered for police use. The Glock 45 frame features a Glock 17 full size grip length with a shorter Glock 19 length dust cover and is fitted with a Glock 19 slide with front slide serrations. Unlike the Glock 19X, the Glock 45 features a flared magazine well and can accommodate Gen 5 magazines, because unlike the 19X, the Glock 45 does not have the front lip that blocks the Gen 5 magazine's extended floor plate. The Glock 45 also deletes the lanyard loop found at the back of the grip on the 19X and comes standard with plastic sights instead of the night sights standard on the 19X.
Glock 46: The Glock 46 is a "compact" version like the Glock 19. This model has a rotating barrel breech lock system. It had been designed as option to bid for a service pistol, with law enforcement agencies in Germany at state and federal level in mind. The differing breech lock system makes it possible to disassemble the firearm without the need to press the trigger, but only with no projectile chambered. Also, enhanced drop-safety is a must – the model complies with the specifications in the German technical guideline (Technische Richtlinie "Pistole") for service pistols. Police in Saxony-Anhalt chose the Glock 46 TR among three competitors and are to receive up to 8,600 new pistols until 2021 for over 6,400 officers, replacing their ageing Pistole 6 (P6), a SIG Sauer P225 variant. The state of Saxony-Anhalt is the first introducing a pistol made by Glock as standard-issue sidearm into its force.
Glock 47: The Glock 47 is a full-sized handgun created for the U.S. Customs and Border Protection, who wanted a version of the Glock 17 that has full parts compatibility with a Glock 19, saving for the slide and barrel. This means that the G47 slide and barrel can be put on a G19 frame to give the G19 a longer slide, barrel, and sight radius (the equivalent of a Glock 49), and the G19 slide and barrel can be put on a G47 frame to create a pistol that functions like a Glock 45. The G47 recoil spring assembly is the same as for the G19. The G47 frame is the same as for the G45 with a shorter dust cover, and the G17-length slide is modified to fill in the gap in a similar manner to the G34. The G47 also comes with MOS cuts and magazines that have the Gen 5 extended base plate. In short, G47 and G19 Gen5/MOS/MOD1/FS have modularity between both pistols. Reportedly, the US Secret Service is also using this model, and the G47 Gen 5 MOS is now available to the public.
Glock 48: The Glock 48 is a "slimline" version similar to the subcompact Glock 43 and 43X. All components of the Glock 48 are identical to the Glock 43X except the slide and barrel which are longer. The slides for the Glock 43, Glock 43X, and Glock 48 are functional on any of those three frames. The G48 features a 4.17-inch-long barrel, front slide serrations, a built-in extended beaver tail, a reversible magazine catch (similar to Gen 5 models), GMB rifling (again similar to Gen 5 models), a two-tone finish (silver slide/black receiver), and a magazine capacity of 10 rounds. In the E.U., Glock offers the 48 and the 43X with a rail. At least three aftermarket sources manufacture 15-round flush-fit magazines for the Glock 43X and Glock 48, which make the Glock 43X and Glock 48 match the standard capacity of the Glock 19 in a narrower pistol.
Glock 49: In May 2022, a photograph was posted by noted firearms author Larry Vickers, who was asking for more information about the firearm pictured in the photograph. The weapon, marked Glock 49, appeared to be a Glock 47 with a shorter Glock 19 size grip and magazine. The pictured Gen 5 firearm lacked any MOS slide cuts. The handgun in that photo was apparently a sample and the Glock 49 has not been marketed in that form. Later, in November 2023, American firearms distributor TALO Distributors announced the introduction of the Glock 49. Like Larry Vicker's photograph, it used the full-length Glock 47 slide and the Glock 19 frame. It differs from Vickers' photograph in that the slide is cut for the MOS optics option. The model is sold exclusively through TALO Distributors, and Glock does not offer to sell it directly. The new model is essentially a crossover similar to the Glock 45, but in reverse. Instead of a compact slide mated to a full-grip frame, it is a full-length slide mated to a compact frame.
9×21mm
The 9×21mm pistol cartridge was adopted and commercialized by Israel Military Industries for those markets (such as in France, Italy, and Mexico) where military service cartridges like the 9×19mm Parabellum are banned for civilian use. Glock produces 9×21mm-chambered versions of the Model 17, and other 9mm Parabellum models, for these markets. These alternative caliber versions are marked with the same model number as the 9mm Parabellum version only differing in the barrel and caliber marking. Glock does not export or produce the 9×21mm pistols for the United States commercial market. This makes any 9×21mm Glock model a unique and are a highly desirable item for US firearm collectors. A limited number of 9×21mm Glocks have found their way into the US and are mostly held by collectors and gun enthusiasts.
10mm Auto
Glock 20: The Glock 20, introduced in 1991, was developed for the then-growing law enforcement and security forces market for the 10mm Auto. The pistol handles both full-power and reduced "FBI" loads that have reduced muzzle velocity. Due to the longer cartridge and higher pressures, the pistol is slightly larger than the Glock 17, having a roughly greater width and greater length. Though many small parts interchange with the Glock 17, with a close to 50% parts commonality, the major assemblies are scaled-up and do not interchange. The standard magazine capacity of the Glock 20 is 15 rounds, and there are no other factory magazine offerings. In 2009, Glock announced they would offer a barrel as a drop-in option.
Glock 20SF: The 20SF is a version of the Glock 20 that uses the Short Frame (SF) which is based on the standard G20 frame (same width), but reduces the trigger reach from the back of the grip by and the heel of the pistol is shortened by so the trigger can be reached and operated better by users with relatively small hands.
Glock 29: The Glock 29 is a 10mm Auto equivalent of the subcompact Glock 26 introduced in 1997 along with the Glock 30 (.45 ACP). The pistol features a barrel and a standard magazine capacity of 10 rounds. Like other subcompact Glock pistols, the Glock 29 functions with the factory magazines from its related full-size model, giving an optional capacity of 15 rounds.
Glock 29SF: The 29SF version of the Glock 29 uses the SF which is based on the standard G29 frame (same width), but reduces the trigger reach from the back of the grip by .
Glock 40: The Glock 40, introduced in 2015, is a 10mm Auto equivalent of the long-slide Glock 17L. The Glock 40 is only made with the "Gen4" frame and "MOS" (Modular Optic System) configuration. The Glock 40 uses the same magazine as the Glock 20.
.45 ACP
Glock pistols chambered for the .45 ACP (and the .45 GAP) feature octagonal polygonal rifling rather than the hexagonal-shaped bores used for models in most other chamberings. Octagonal rifling provides a better gas seal in relatively large diameter rifled bores, since an octagon resembles a circle more closely than a hexagon.
Glock 21: The Glock 21 is a .45 ACP version of the Glock 20 designed primarily for the American market. Compared to the Glock 20 chambered in 10mm Auto, the slide of the Glock 21 is lighter to compensate for the lower-energy .45 ACP cartridge. The standard Glock 21 magazine is of the single-position-feed, staggered-column type with a capacity of 13 rounds. While there is no other factory magazine, a "+1" base plate is available to make the magazine 14-round.
Glock 21SF: The 21SF is a version of the Glock 21 that uses a Short Frame lower which is based on the standard G21 frame (same width), but reduces trigger reach from the back of the grip by , and the heel of the pistol is shortened by so the trigger can be reached and operated better by users with smaller hands.
Glock 30: The Glock 30 is a .45 ACP version of the subcompact Glock 29, with a standard magazine capacity of 10 rounds. The standard magazine includes a "+1" base plate, and can be made into a 9-round magazine by swapping in a flush base plate. The factory magazine from the Glock 21, with a capacity of 13 rounds (14 rounds with the +1 base plate), will function in the Glock 30.
Glock 30SF: The 30SF is a version of the Glock 30 that uses a Short Frame lower which is based on the standard G30 frame (same width), but reduces trigger reach from the back of the grip by . The G30SF accepts the same double-stack .45ACP magazines as the G30 and G21.
Glock 30S: The 30S is a version of the Glock 30 that features a thin slide (same slide as the G36), a Short Frame lower, and the same double stack magazines as the Glock 30 and 21. Like the G30, G30S magazines holds 10 rounds.
Glock 36: The Glock 36 is a "slimline" version of the subcompact Glock 30 that features an ultracompact slide and frame and is chambered for the .45 ACP cartridge. The Glock 36 is the first Glock pistol to be manufactured with a single-stack magazine, having a standard capacity of six rounds and being unique to the model. The Glock 36 cannot use factory magazines from its larger relatives due to its single-stack magazine design.
Glock 41: The Glock 41 is a competition version of the Glock 21, much like what the G34 is in relation to the G17; it features a 5.3-inch barrel and an elongated slide. The Glock 41 is only made with the "Gen4" frame.
.40 S&W
In 1990, Smith & Wesson and Winchester developed the .40 Smith & Wesson by shortening the 10mm case. This created a round that was more powerful than the 9mm Parabellum but with more manageable recoil. The round was also still operable in the smaller frame size used for Glock 9mm models. As is typical of many pistols chambered in .40 S&W, each of the standard Glock models (22, 23, and 27) may be easily converted to the corresponding .357 SIG chambering (Glock 31, 32, and 33, respectively) simply by replacing the barrel. No other parts need to be replaced, as the .40 S&W magazines will feed the .357 SIG rounds.
Glock 22: The Glock 22 is a .40 S&W version of the full-sized Glock 17 introduced in 1990. The pistol uses a modified slide, frame, and barrel to account for the differences in size and power of the .40 S&W cartridge. The standard magazine capacity is 15 rounds. A 10-round version of this magazine is also offered for those markets where the magazine capacity for handguns is restricted. Additionally, the standard magazine is available with a "+1" base plate to make it a 16-round magazine. And, there is a Glock factory 22-round extended magazine offered. One can also swap in the +1 base plate to make that a 23-round magazine. The Glock Model 22 is favored and used by multiple law enforcement agencies around the world, including the Baltimore Police Department, Los Angeles Police Department, Miami Police Department, Maryland State Police, Cumberland County Sheriff's Department (NJ), Overland Park Police Department, Kansas City Police Department, Missouri State Highway Patrol, and Alaska State Troopers in the United States; the NSW Police Force, Queensland Police Service and the Northern Territory Police Force in Australia; the Edmonton Police Service, Calgary Police Service, Alberta Sheriffs Branch, Winnipeg Police Service, Toronto Police Service, Ottawa Police Service, and British Columbia Sheriff Service in Canada; and the National Bureau of Investigation (Philippines).
Glock 23: The Glock 23 is a .40 S&W version of the compact Glock 19. It is dimensionally identical to the Glock 19, but is slightly heavier and uses a modified slide, frame, .40 S&W barrel, and a standard magazine capacity of 13 rounds. For jurisdictions limiting the magazine capacity of handguns, Glock offers a 10-round version of the standard magazine as well. There is also a factory 14-round version of the standard magazine using the +1 base plate. Finally, any magazine made for the Glock 22 will work in the Glock 23.
Glock 24: The Glock 24 is a .40 S&W long-slide variant of the Glock 22, similar in concept to the Glock 17L. Additionally, a compensated, ported-barrel version designated the 24C was also produced. The Glock 24 was introduced in 1994 and officially dropped from the company's regular product lineup upon the release of the Glock 34 and 35. The Glock 24 can use any magazine made for the Glock 22.
Glock 27: The Glock 27 is a .40 S&W version of the subcompact Glock 26, with a standard magazine capacity of 9 rounds. Glock also offers a 10-round version of this magazine with the +1 base plate. The factory magazines from the larger Glock 22 and 23 will function in the Glock 27, increasing capacity to 13, 14, 15, 16 or 22 rounds. Spacers are available that fit on these larger-capacity magazines themselves; they have the effect of "extending" the magazine well of the pistol, thereby improving the ergonomic feel of the pistol when the longer magazines are inserted.
Glock 35: The Glock 35 is a .40 S&W version of the competition Glock 34. The Glock Model 35 was the service pistol for the Kentucky State Police, but by the summer of 2017, they had reverted from the Glock 35 back to 9mm weapons because of penetration improvements in the 9mm bullets, noting some officers had never been able to make the switch from 9mm to .40 S&W in the first place due to their struggles in mastering the higher caliber. The Glock 35 can use any magazine made for the Glock 22.
.380 ACP
The first two .380 ACP models (Glock 25 and 28) were released in 1995 to provide a less powerful alternative to the 9mm Parabellum and 9×21mm, and primarily intended for markets (such as Brazil) that prohibit civilian ownership of firearms chambered in more powerful calibers. Made in Austria, they are legally prohibited from being imported for civilians in the United States, lacking sufficient points to meet the import restrictions. Recently, a limited run of the Glock 28 were manufactured in Glock's US plant for sale in the US. The Glock 25 and 28 are also prohibited from ownership in Canada due to not meeting minimum barrel length requirements for handguns.
Due to the relatively low bolt thrust of the .380 ACP cartridge, the locked-breech design of the Glock 19 and Glock 26 was minimally modified for the Glock 25 and Glock 28 to implement unlocked breech operation. It operates via straight blowback of the slide. This required modification of the locking surfaces on the barrel, as well as a redesign of the former locking block. Unusual for a blowback design, the barrel is not fixed to the frame. It moves rearward in recoil until it is tilted below the slide, similar to the standard locked-breech system. The reduced size and mass of the Glock 42 allowed for the return to the Glock-standard locked-breech design.
Glock 25: The Glock 25 , introduced in 1995, is a blowback derivative of the compact ( barrel) Glock 19. The magazine capacity is 15 rounds. Standard fixed sight elevation is 6.9 mm, unlike the 6.5 mm elevation used for the 9×19mm models.
Glock 28: The Glock 28 , introduced in 1997, is a blowback derivative of the subcompact ( barrel) Glock 26. The standard magazine capacity is 10 rounds, but the 15-round Glock 25 magazine will function in the Glock28. Standard fixed-sight elevation is 6.9 mm, unlike the 6.5 mm elevation used for the 9×19mm Parabellum models.
Glock 42: The Glock 42 , introduced in 2014, is a locked-breech "slimline" ( barrel) design. The single-stack magazine has a capacity of six rounds. It is Glock's smallest model ever made and is manufactured in the United States, which unlike the Glock25 and 28, allows it to be sold in the domestic USA market.
.357 SIG
In 1994, SIG and Federal developed the .357 SIG to match the ballistics of typical .357 Magnum loads. This was done by necking down a shortened 10mm case to .357, so it was easy for Glock to alter their .40 Smith & Wesson models to the new caliber. As is typical of pistols chambered in .357 SIG, each of the standard Glock models (31, 32, and 33) may be easily converted to the corresponding .40 S&W chambering (Glock 22, 23, and 27, respectively) simply by replacing the barrel. No other parts need to be replaced, as the .357 SIG magazines will feed the .40 S&W round. Though marked as a different caliber, the .357 SIG magazines are the same as for the .40 S&W models and interchangeable.
Glock 31: The Glock 31 is a .357 SIG variant of the full-sized Glock 22. The standard magazine capacity of the Glock 31 is 15 rounds. It can accept magazines intended for the Glock 22 as well.
Glock 32: The Glock 32 is a .357 SIG variant of the compact Glock 23. The standard magazine capacity of the Glock 32 is 13 rounds. It can accept magazines intended for the Glock 31, Glock 22, or Glock 23 as well.
Glock 33: The Glock 33 is a .357 SIG variant of the subcompact Glock 27. The standard magazine capacity of the Glock 33 is 9 rounds. It can also accept magazines intended for the Glock 32, Glock 31, Glock 22, Glock 23, and Glock 27.
.45 GAP
Glock pistols chambered for the .45 GAP (and the .45 ACP) feature octagonal polygonal rifling rather than the hexagonal-shaped bores used for models in most other chamberings. Octagonal rifling provides a better gas seal in relatively large diameter rifled bores, since an octagon will have shorter sides and shallower angles than a hexagon.
Glock 37: The Glock 37 is a .45 GAP version of the Glock 17. It uses a wider, beveled slide, larger barrel, and different magazine, but is otherwise similar to the Glock 17. The Glock 37 first appeared in 2003. It was designed to offer ballistic performance comparable with the .45 ACP in the frame size of the Glock 17. The concern with the size of the Glock 20/21 has been addressed by the Glock 36, 21SF, and 30SF, all of which featured reduced-size frames. The standard magazine capacity of the Glock 37 is 10 rounds.
Glock 38: The Glock 38 is a .45 GAP version of the compact Glock 19. The standard magazine capacity of the Glock 38 is 8 rounds, but it can use the 10-round magazines of the Glock 37.
Glock 39: The Glock 39 is a .45 GAP version of the subcompact Glock 26. The standard magazine capacity of the Glock 39 is 6 rounds, but it can use the 8- and 10-round magazines of the Glock 37 and Glock 38.
.22 Long Rifle
Glock 44: The Glock 44 is a .22 Long Rifle rimfire model based on the Glock 19. While the Glock 44 is similar in size as the Glock 19, the Glock 44 has a magazine capacity of 10 rounds and uses a simple blowback mechanism instead of a locked breech mechanism used on nearly all other Glock pistols. In addition, the Glock 44 is lighter than the G19, weighing in at just 12 ounces, and it retains Glock's popular polygonal rifling, which has been tested by Glock to work with lead .22 bullets. It uses a steel/polymer composite slide due to the lower slide mass required to function with the less powerful cartridge. The 44 features the company's proprietary "Safe Action" trigger and a footprint so exact that holsters made for the 19 or 23 are interchangeable with the 44.
Production in other countries
: Russian firms such as Skat, ORSIS, and Izhmash assemble three models of Glock pistols locally: the Glock 17, 34, and 35.
: As of 2015, there were plans to assemble Glock 17 pistols at army workshops in Uruguay to fulfill the needs of the national military services and law enforcement organizations.
: Glock pistols are manufactured by the Glock Inc. subsidiary division located in the United States. Those batches are identical to the Austrian-made ones, but they are marked as "USA", instead of "AUSTRIA", on the slide.
Clones
Third party frames and slides for Glock pistols began to appear in the early 2000s. This has led to "Glock" becoming a generic term encompassing pistols not made by Glock Ges.m.b.H. especially as expiring patents allowed for complete Glock clones to be made. As of 2019, a large number of American companies produced Glock clones.
: In 2017, it was reported that Norinco was able to make a clone of the Glock 17 known as the NP-7 (or NP7). The pistol was subcontracted to Hunan Ordnance Industry Group through the Hunan Ordnance and Light Weapons Research Institute. Its features appear to be influenced by the fourth gen Glock 17. The NP7 is being marketed for export sales.
: There are three sidearms made by Iranian DIO's Shahid Kaveh Industry Complex which they call Ra'ad (has a safety selector, possibly an unlicensed copy of Glock 17), Glock 19 and Kaveh-17 (probably an improved Ra'ad, a variant of Glock 17S), all of which are unlicensed clones of Glock pistols. It is not known if they could be adopted by the Iranian military and replace the Browning Hi-Power, M1911 and SIG P226 pistols and they were possibly some prototypes and have never gone on mass production.
: The Tatmadaw have adopted a clone of the Glock 17 known as the MA5 MK II, which was first reported in 2018. They are currently being manufactured and adopted for Myanmar's special forces units. The pistol has a rounded trigger guard and the Tatmadaw emblem on the grip, the latter having rectangular spacing instead of the smaller dot-like spacing.
: Unlicensed Glock clones are made in Pakistan's Khyber region, which were first reported in 2018.
: The 205th Arsenal in Taiwan produces a copy of the Glock 19, named the T97 pistol. The Taiwan-made Glocks were made to replace the Smith & Wesson Model 5906 used by the Taiwan police, but it ultimately did not enter service.
: Akdal Arms produces a pistol named the Ghost TR01, which is heavily influenced by Glock pistols in its design.
: At the Vietnam Defence Expo 2024 in December 2024, Z111 Factory unveiled a new series of Glock clone pistols known as the SN19 (Glock 17 Gen 5), the SN19-T (Steel Frame version) and the SN7VN-M24 chambered in 7.62x25 Tokarev.
Users
Criminal use
Glock pistols have been used in mass shootings in the United States including the 1991 Luby's shooting, the 2007 Virginia Tech shooting, the 2011 Tucson shooting, the 2012 Aurora shooting, the 2012 Sandy Hook Elementary School shooting, the 2015 Charleston church shooting, the 2016 Orlando nightclub shooting, the Pittsburgh synagogue shooting, and the 2022 New York City Subway attack. Glock pistols have also been used in mass shootings elsewhere in the world, including the 2001 Nepalese royal massacre, the 2002 Erfurt massacre, the killings committed by Viktor Kalivoda in 2005, the 2011 Norway attacks, the 2016 Munich shooting, and the 2023 Prague shooting. In April 2022, Ilene Steur, a survivor of the 2022 NYC subway attack, sued Glock and its Austrian parent company for compensation for her physical injuries and emotional pain.
A 2014 report by the Chicago Police Department found that Glock pistols were the third most traced handgun, coming after those from Smith & Wesson and Sturm, Ruger & Co. Experts on gun control, mass shootings, and defense training have cited factors such as reliability, ease of use, and commonness for why Glock pistols are so often involved in mass shootings and other criminal acts. The criminal use of handguns including Glocks has led to calls for increased gun control in the United States. This common usage, however, has been described by Paul M. Barrett to be a result of Glock's overall popularity and market presence in the US and that "this level of violence isn't necessarily tied to a particular[,] to a brand".
In the late 80s, gun control advocates had focused on Glock pistols because of their magazine capacity (compared to six-shot revolvers), but also their "futuristic, distinct appearance".
Glock pistols were singled out for restriction by some jurisdictions and were branded the "hijacker's special" based on the false assumption that they could bypass airport metal detectors because of their polymer frame. This was refuted in Congressional hearings by the ATF, FAA, and other organizations responsible for airline security, which proved embarrassing for the bans' advocates and provided significant publicity for Glock.
In December 2024, the states of Minnesota and New Jersey sued Glock over the design of the 9-millimeter semiautomatic pistol, claiming that the company had failed to make changes to prevent the easy conversion of the pistol into an illegal and much more dangerous machine gun. The conversion is done using a small, illegal device known as a Glock switch.
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3084010 | https://en.wikipedia.org/wiki/Gatehouse | Gatehouse | A gatehouse is a type of fortified gateway, an entry control point building, enclosing or accompanying a gateway for a town, religious house, castle, manor house, or other fortification building of importance. Gatehouses are typically the most heavily armed section of a fortification, to compensate for being structurally the weakest and the most probable attack point by an enemy. There are numerous surviving examples in France, Austria, Germany, England and Japan.
History
Gatehouses made their first appearance in the early antiquity when it became necessary to protect the main entrance to a castle or town. Famous early examples of such gates are those such as the Ishtar Gate in Babylon. Over time, they evolved into very complicated structures with many lines of defence. The Romans began building fortified walls and structures throughout Europe such as the Aurelian Walls of Rome with gates such as Porta San Paolo and Porta Nigra from the ancient defenses of Trier in Germany. Strongly fortified gatehouses would normally include a drawbridge, one or more portcullises, machicolations, arrow loops and possibly even murder-holes where stones would be dropped on attackers. In some castles, the gatehouse was so strongly fortified it took on the function of a keep, sometimes referred to as a "gate keep". In the late Middle Ages, some of these arrow loops might have been converted into gun loops (or gun ports).
Urban defences would sometimes incorporate gatehouses such as Monnow Bridge in Monmouth. York has four important gatehouses, known as "Bars", in its city walls including the Micklegate Bar.
The French term for gatehouse is logis-porche. This could be a large, complex structure that served both as a gateway and lodging or it could have been composed of a gateway through an enclosing wall. A very large gatehouse might be called a châtelet (small castle).
At the end of the Middle Ages, many gatehouses in England and France were converted into beautiful, grand entrance structures to manor houses or estates. Many of them became a separate feature free-standing or attached to the manor or mansion only by an enclosing wall. By this time the gatehouse had lost its defensive purpose and had become more of a monumental structure designed to harmonise with the manor or mansion.
England
Bargate, in Hampshire is a medieval gatehouse in the city centre of Southampton, England. Constructed in 1180 as part of the Southampton town walls
Ightham Mote, in Kent has an imposing 13th and 14th century gatehouse.
Durham Castle, in Durham has an 11th-century gatehouse that is now used as accommodation for students attending University College, Durham.
Layer Marney Tower, the epitome of the Tudor gatehouse.
Stokesay Castle, a 13th-century fortified manor house in Shropshire has a Jacobean half-timbered gatehouse.
Stanway House, Stanway, Gloucestershire, where the gatehouse measures 44 ft. by 22 ft. and has three storeys.
Westwood House, Worcestershire, which has a frontage of 54 ft. with two storeys.
Burton Agnes Hall, East Riding of Yorkshire, which has three storeys and is flanked by great octagonal towers at the angles.
Hylton Castle, Hylton, Sunderland, although it is an actual castle, it is styled in the shape of a classical gatehouse (this is due to the castle being built for comfort as opposed to a castle for defence).
France
Château de Châteaubriant, two gatehouses (13th and 14th century), one for the lower bailey, one for the upper ward.
Château de Suscinio, a large 15th-century gatehouse in the logis-porte style, Morbihan, Brittany.
Château de Trécesson, a simple 14th-century gatehouse on a moated manor house in Morbihan, Brittany (see French Wikipedia page, Château de Trécesson)
Château de Vitré, a large 15th-century châtelet or gatehouse in Ille-et-Vilaine, Brittany (see French Wikipedia page, Château de Vitré)
United States
Latrobe Gate, a Greek Revival and Italianate gatehouse built in 1806, Washington, D.C.
Lorraine Park Cemetery Gate Lodge, a Queen Anne style stone and frame building constructed in 1884, Woodlawn, Baltimore County, Maryland.
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3086453 | https://en.wikipedia.org/wiki/Core%20sample | Core sample | A core sample is a cylindrical section of (usually) a naturally-occurring substance. Most core samples are obtained by drilling with special drills into the substance, such as sediment or rock, with a hollow steel tube, called a core drill. The hole made for the core sample is called the "core hole". A variety of core samplers exist to sample different media under different conditions; there is continuing development in the technology. In the coring process, the sample is pushed more or less intact into the tube. Removed from the tube in the laboratory, it is inspected and analyzed by different techniques and equipment depending on the type of data desired.
Core samples can be taken to test the properties of manmade materials, such as concrete, ceramics, some metals and alloys, especially the softer ones. Core samples can also be taken of living things, including human beings, especially of a person's bones for microscopic examination to help diagnose diseases.
Methods
The composition of the subject materials can vary from almost liquid to the strongest materials found in nature or technology, and the location of the subject materials can vary from on the laboratory bench to over 10 km from the surface of the Earth in a borehole. The range of equipment and techniques applied to the task is correspondingly great. Core samples are most often taken with their long axis oriented roughly parallel to the axis of a borehole, or parallel to the gravity field for the gravity-driven tools. However it is also possible to take core samples from the wall of an existing borehole. Taking samples from an exposure, be it an overhanging rock face or on a different planet, is almost trivial. (The Mars Exploration Rovers carry a Rock Abrasion Tool, which is logically equivalent to the "rotary sidewall core" tool described below.)
Some common techniques include:
gravity coring, in which the core sampler is dropped into the sample, usually the bed of a water body, but essentially the same technique can also be done on soft materials on land. The penetration forces, if recorded, give information about the strength of different depths in the material, which may be the only information required, with samples as an incidental benefit. This technique is common in both civil engineering site investigations (where the technique shades into pile driving) and geological studies of recent aquatic deposits. The low strength of the materials penetrated means that cores have to be relatively small.
vibrating, in which the sampler is vibrated to allow penetration into thixotropic media. Again, the physical strength of the subject material limits the size of the core that can be retrieved.
drilling exploration diamond drilling where a rotating annular tool backed up by a cylindrical core sample storage device is pressed against the subject materials to cut out a cylinder of the subject material. A mechanism is normally needed to retain the cylindrical sample in the coring tool. Depending on circumstances, particularly the consistency and composition of the subject materials, different arrangements may be needed within the core tools to support and protect the sample on its way to the surface; it is often also necessary to control or reduce the contact between the drilling fluid and the core sample, to reduce changes from the coring process. The mechanical forces imposed on the core sample by the tool frequently lead to fracture of the core and loss of less-competent intervals, which can greatly complicate the interpretation of the core. Cores can routinely be cut as small as a few millimeters in diameter (in wood, for dendrochronology) up to over 150 millimeters in diameter (routine in oil exploration). The lengths of samples can range from less than a meter (again, in wood, for dendrochronology) up to around 200 meters in one run, though 27 to 54 meters is more usual (in oil exploration), and many runs can be made in succession if "quick look" analysis in the field suggests that the zone of interest is continuing.
percussion sidewall coring uses robust cylindrical "bullets" explosively propelled into the wall of a borehole to retrieve a (relatively) small, short core sample. These tend to be heavily shattered, rendering porosity/ permeability measurements dubious, but are often sufficient for lithological and micropalaeontological study. Many samples can be attempted in a single run of the tools, which are typically configured with 20 to 30 "bullets" and propulsive charges along the length of a tool. Several tools can often be ganged together for a single run. The success rates for firing a particular bullet, it penetrating the borehole wall, the retention system recovering the bullet from the borehole wall, and the sample is retained in the bullet are all relatively low, so it is not uncommon for only half the samples attempted to be successful. This is an important consideration in planning sample programs.
rotary sidewall coring where a miniaturized automated rotary drilling tool is applied to the side of the borehole to cut a sample similar in size to a percussion sidewall core (described above). These tend to suffer less deformation than percussion cores. However, the core-cutting process takes longer and jams are common in the ancillary equipment which retrieves the sample from the drill bit and stores it within the tool body.
coring by hand can be done with a variety of instruments, such as a Russian Peat Corer, a handheld piston corer or simply a hollow tube. The advantages of coring by hand are low costs and quick operation, but a handheld corer will only reach limited depths, ranging from a few decimeters in clay soils to a few meters in soft lake sediments.
Management of cores and data
Although often neglected, core samples always degrade to some degree in the process of cutting the core, handling it, and studying it. Non-destructive techniques are increasingly common, e.g., the use of MRI scanning to characterize grains, pore fluids, pore spaces (porosity) and their interactions (constituting part of permeability) but such expensive subtlety is likely wasted on a core that has been shaken on an unsprung lorry for 300 km of dirt road. What happens to cores between the retrieval equipment and the final laboratory (or archive) is an often neglected part of record keeping and core management.
Coring has come to be recognized as an important source of data, and more attention and care is being put on preventing damage to the core during various stages of it transportation and analysis. The usual way to do this is to freeze the core completely using liquid nitrogen, which is cheaply sourced. In some cases, special polymers are also used to preserve and seat/cushion the core from damage.
Equally, a core sample which cannot be related to its context (where it was before it became a core sample) has lost much of its benefit. The identification of the borehole, and the position and orientation ("way up") of the core in the borehole is critical, even if the borehole is in a tree trunk – dendrochronologists always try to include a bark surface in their samples so that the date of most-recent growth of the tree can be unambiguously determined.
If these data become separated from core samples, it is generally impossible to regain that data. The cost of a coring operation can vary from a few currency units (for a hand-caught core from a soft soil section) to tens of millions of currency units (for sidewall cores from a remote-area offshore borehole many kilometres deep). Inadequate recording of such basic data has ruined the utility of both types of core.
Different disciplines have different local conventions of recording these data, and the user should familiarize themselves with their area's conventions. For example, in the oil industry, orientation of the core is typically recorded by marking the core with two longitudinal colour streaks, with the red one on the right when the core is being retrieved and marked at surface. Cores cut for mineral mining may have their own, different, conventions. Civil engineering or soil studies may have their own, different, conventions as their materials are often not competent enough to make permanent marks on.
It is becoming increasingly common to retain core samples in cylindrical packaging which forms part of the core-cutting equipment, and to make the marks of record on these "inner barrels" in the field prior to further processing and analysis in the laboratory. Sometimes core is shipped from the field to the laboratory in as long a length as it comes out of the ground; other times it is cut into standard lengths (5m or 1m or 3 ft) for shipping, then reassembled in the laboratory. Some of the "inner barrel" systems are capable of being reversed on the core sample, so that in the laboratory the sample goes "wrong way up" when the core is reassembled. This can complicate interpretation.
If the borehole has petrophysical measurements made of the wall rocks, and these measurements are repeated along the length of the core then the two data sets correlated, one will almost universally find that the depth "of record" for a particular piece of core differs between the two methods of measurement. Which set of measurements to believe then becomes a matter of policy for the client (in an industrial setting) or of great controversy (in a context without an overriding authority). Recording that there are discrepancies, for whatever reason, retains the possibility of correcting an incorrect decision at a later date ; destroying the "incorrect" depth data makes it impossible to correct a mistake later. Any system for retaining and archiving data and core samples needs to be designed so that dissenting opinion like this can be retained.
If core samples from a campaign are competent, it is common practice to "slab" them – cut the sample into two or more samples longitudinally – quite early in laboratory processing so that one set of samples can be archived early in the analysis sequence as a protection against errors in processing. "Slabbing" the core into a 2/3 and a 1/3 set is common. It is also common for one set to be retained by the main customer while the second set goes to the government (who often impose a condition for such donation as a condition of exploration/ exploitation licensing). "Slabbing" also has the benefit of preparing a flat, smooth surface for examination and testing of profile permeability, which is very much easier to work with than the typically rough, curved surface of core samples when they're fresh from the coring equipment. Photography of raw and "slabbed" core surfaces is routine, often under both natural and ultra-violet light.
A unit of length occasionally used in the literature on seabed cores is cmbsf, an abbreviation for centimeters below sea floor.
History of coring
The technique of coring long predates attempts to drill into the Earth’s mantle by the Deep Sea Drilling Program. The value to oceanic and other geologic history of obtaining cores over a wide area of sea floors soon became apparent. Core sampling by many scientific and exploratory organizations expanded rapidly. To date hundreds of thousands of core samples have been collected from floors of all the planet's oceans and many of its inland waters.
Access to many of these samples is facilitated by the Index to Marine & Lacustrine Geological Samples.
Informational value of core samples
Coring began as a method of sampling surroundings of ore deposits and oil exploration. It soon expanded to oceans, lakes, ice, mud, soil and wood. Cores on very old trees give information about their growth rings without destroying the tree.
Cores indicate variations of climate, species and sedimentary composition during geologic history. The dynamic phenomena of the Earth's surface are for the most part cyclical in a number of ways, especially temperature and rainfall.
There are many ways to date a core. Once dated, it gives valuable information about changes of climate and terrain. For example, cores in the ocean floor, soil and ice have altered the view of the geologic history of the Pleistocene entirely.
Alternatives
Reverse circulation drilling is a method in which rock cuttings are continuously extracted through the hollow drill rod and can be sampled for analysis. The method may be faster and use less water than core drilling, but does not produce cores of relatively undisturbed material, so less information on the rock structure can be derived from analysis. If compressed air is used for cutting extraction the sample remains uncontaminated, is available almost immediately, and the method has a low environmental impact.
| Physical sciences | Geology: General | Earth science |
3087602 | https://en.wikipedia.org/wiki/Topological%20order | Topological order | In physics, topological order is a kind of order in the zero-temperature phase of matter (also known as quantum matter). Macroscopically, topological order is defined and described by robust ground state degeneracy and quantized non-abelian geometric phases of degenerate ground states. Microscopically, topological orders correspond to patterns of long-range quantum entanglement. States with different topological orders (or different patterns of long range entanglements) cannot change into each other without a phase transition.
Various topologically ordered states have interesting properties, such as (1) topological degeneracy and fractional statistics or non-abelian group statistics that can be used to realize a topological quantum computer; (2) perfect conducting edge states that may have important device applications; (3) emergent gauge field and Fermi statistics that suggest a quantum information origin of elementary particles; (4) topological entanglement entropy that reveals the entanglement origin of topological order, etc. Topological order is important in the study of several physical systems such as spin liquids, and the quantum Hall effect, along with potential applications to fault-tolerant quantum computation.
Topological insulators and topological superconductors (beyond 1D) do not have topological order as defined above, their entanglements being only short-ranged, but are examples of symmetry-protected topological order.
Background
Matter composed of atoms can have different properties and appear in different forms, such as solid, liquid, superfluid, etc. These various forms of matter are often called states of matter or phases. According to condensed matter physics and the principle of emergence, the different properties of materials generally arise from the different ways in which the atoms are organized in the materials. Those different organizations of the atoms (or other particles) are formally called the orders in the materials.
Atoms can organize in many ways which lead to many different orders and many different types of materials. Landau symmetry-breaking theory provides a general understanding of these different orders. It points out that different orders really correspond to different symmetries in the organizations of the constituent atoms. As a material changes from one order to another order (i.e., as the material undergoes a phase transition), what happens is that the symmetry of the organization of the atoms changes.
For example, atoms have a random distribution in a liquid, so a liquid remains the same as we displace atoms by an arbitrary distance. We say that a liquid has a continuous translation symmetry. After a phase transition, a liquid can turn into a crystal. In a crystal, atoms organize into a regular array (a lattice). A lattice remains unchanged only when we displace it by a particular distance (integer times a lattice constant), so a crystal has only discrete translation symmetry. The phase transition between a liquid and a crystal is a transition that reduces the continuous translation symmetry of the liquid to the discrete symmetry of the crystal. Such a change in symmetry is called symmetry breaking. The essence of the difference between liquids and crystals is therefore that the organizations of atoms have different symmetries in the two phases.
Landau symmetry-breaking theory has been a very successful theory. For a long time, physicists believed that Landau Theory described all possible orders in materials, and all possible (continuous) phase transitions.
Discovery and characterization
However, since the late 1980s, it has become gradually apparent that Landau symmetry-breaking theory may not describe all possible orders. In an attempt to explain high temperature superconductivity the chiral spin state was introduced. At first, physicists still wanted to use Landau symmetry-breaking theory to describe the chiral spin state. They identified the chiral spin state as a state that breaks the time reversal and parity symmetries, but not the spin rotation symmetry. This should be the end of the story according to Landau's symmetry breaking description of orders. However, it was quickly realized that there are many different chiral spin states that have exactly the same symmetry, so symmetry alone was not enough to characterize different chiral spin states. This means that the chiral spin states contain a new kind of order that is beyond the usual symmetry description. The proposed, new kind of order was named "topological order". The name "topological order" is motivated by the low energy effective theory of the chiral spin states which is a topological quantum field theory (TQFT). New quantum numbers, such as ground state degeneracy (which can be defined on a closed space or an open space with gapped boundaries, including both Abelian topological orders and non-Abelian topological orders) and the non-Abelian geometric phase of degenerate ground states, were introduced to characterize and define the different topological orders in chiral spin states. More recently, it was shown that topological orders can also be characterized by topological entropy.
But experiments soon indicated that chiral spin states do not describe high-temperature superconductors, and the theory of topological order became a theory with no experimental realization. However, the similarity between chiral spin states and quantum Hall states allows one to use the theory of topological order to describe different quantum Hall states. Just like chiral spin states, different quantum Hall states all have the same symmetry and are outside the Landau symmetry-breaking description. One finds that the different orders in different quantum Hall states can indeed be described by topological orders, so the topological order does have experimental realizations.
The fractional quantum Hall (FQH) state was discovered in 1982 before the introduction of the concept of topological order in 1989. But the FQH state is not the first experimentally discovered topologically ordered state. The superconductor, discovered in 1911, is the first experimentally discovered topologically ordered state; it has Z2 topological order.
Although topologically ordered states usually appear in strongly interacting boson/fermion systems, a simple kind of topological order can also appear in free fermion systems. This kind of topological order corresponds to integral quantum Hall state, which can be characterized by the Chern number of the filled energy band if we consider integer quantum Hall state on a lattice. Theoretical calculations have proposed that such Chern numbers can be measured for a free fermion system experimentally.
It is also well known that such a Chern number can be measured (maybe indirectly) by edge states.
The most important characterization of topological orders would be the underlying fractionalized excitations (such as anyons) and their fusion statistics and braiding statistics (which can go beyond the quantum statistics of bosons or fermions). Current research works show that the loop and string like excitations exist for topological orders in the 3+1 dimensional spacetime, and their multi-loop/string-braiding statistics are the crucial signatures for identifying 3+1 dimensional topological orders. The multi-loop/string-braiding statistics of 3+1 dimensional topological orders can be captured by the link invariants of particular topological quantum field theory in 4 spacetime dimensions.
Mechanism
A large class of 2+1D topological orders is realized through a mechanism called string-net condensation. This class of topological orders can have a gapped edge and are classified by unitary fusion category (or monoidal category) theory. One finds that string-net condensation can generate infinitely many different types of topological orders, which may indicate that there are many different new types of materials remaining to be discovered.
The collective motions of condensed strings give rise to excitations above the string-net condensed states. Those excitations turn out to be gauge bosons. The ends of strings are defects which correspond to another type of excitations. Those excitations are the gauge charges and can carry Fermi or fractional statistics.
The condensations of other extended objects such as "membranes", "brane-nets", and fractals also lead to topologically ordered phases and "quantum glassiness".
Mathematical formulation
We know that group theory is the mathematical foundation of symmetry-breaking orders. What is the mathematical foundation of topological order?
It was found that a subclass of 2+1D topological orders—Abelian topological orders—can be classified by a K-matrix approach. The string-net condensation suggests that tensor category (such as fusion category or monoidal category) is part of the mathematical foundation of topological order in 2+1D. The more recent researches suggest that
(up to invertible topological orders that have no fractionalized excitations):
2+1D bosonic topological orders are classified by unitary modular tensor categories.
2+1D bosonic topological orders with symmetry G are classified by G-crossed tensor categories.
2+1D bosonic/fermionic topological orders with symmetry G are classified by unitary braided fusion categories over symmetric fusion category, that has modular extensions. The symmetric fusion category Rep(G) for bosonic systems and sRep(G) for fermionic systems.
Topological order in higher dimensions may be related to n-Category theory. Quantum operator algebra is a very important mathematical tool in studying topological orders.
Some also suggest that topological order is mathematically described by extended quantum symmetry.
Applications
The materials described by Landau symmetry-breaking theory have had a substantial impact on technology. For example, ferromagnetic materials that break spin rotation symmetry can be used as the media of digital information storage. A hard drive made of ferromagnetic materials can store gigabytes of information. Liquid crystals that break the rotational symmetry of molecules find wide application in display technology. Crystals that break translation symmetry lead to well defined electronic bands which in turn allow us to make semiconducting devices such as transistors. Different types of topological orders are even richer than different types of symmetry-breaking orders. This suggests their potential for exciting, novel applications.
One theorized application would be to use topologically ordered states as media for quantum computing in a technique known as topological quantum computing. A topologically ordered state is a state with complicated non-local quantum entanglement. The non-locality means that the quantum entanglement in a topologically ordered state is distributed among many different particles. As a result, the pattern of quantum entanglements cannot be destroyed by local perturbations. This significantly reduces the effect of decoherence. This suggests that if we use different quantum entanglements in a topologically ordered state to encode quantum information, the information may last much longer. The quantum information encoded by the topological quantum entanglements can also be manipulated by dragging the topological defects around each other. This process may provide a physical apparatus for performing quantum computations. Therefore, topologically ordered states may provide natural media for both quantum memory and quantum computation. Such realizations of quantum memory and quantum computation may potentially be made fault tolerant.
Topologically ordered states in general have a special property that they contain non-trivial boundary states. In many cases, those boundary states become perfect conducting channel that can conduct electricity without generating heat. This can be another potential application of topological order in electronic devices.
Similarly to topological order, topological insulators also have gapless boundary states. The boundary states of topological insulators play a key role in the detection and the application of topological insulators.
This observation naturally leads to a question:
are topological insulators examples of topologically ordered states?
In fact topological insulators are different from topologically ordered states defined in this article.
Topological insulators only have short-ranged entanglements and have no topological order, while the topological order defined in this article is a pattern of long-range entanglement. Topological order is robust against any perturbations. It has emergent gauge theory, emergent fractional charge and fractional statistics. In contrast, topological insulators are robust only against perturbations that respect time-reversal and U(1) symmetries. Their quasi-particle excitations have no fractional charge and fractional statistics. Strictly speaking, topological insulator is an example of symmetry-protected topological (SPT) order, where the first example of SPT order is the Haldane phase of spin-1 chain. But the Haldane phase of spin-2 chain has no SPT order.
Potential impact
Landau symmetry-breaking theory is a cornerstone of condensed matter physics. It is used to define the territory of condensed matter research. The existence of topological order appears to indicate that nature is much richer than Landau symmetry-breaking theory has so far indicated. So topological order opens up a new direction in condensed matter physics—a new direction of highly entangled quantum matter. We realize that quantum phases of matter (i.e. the zero-temperature phases of matter) can be divided into two classes: long range entangled states and short range entangled states. Topological order is the notion that describes the long range entangled states: topological order = pattern of long range entanglements. Short range entangled states are trivial in the sense that they all belong to one phase. However, in the presence of symmetry, even short range entangled states are nontrivial and can belong to different phases. Those phases are said to contain SPT order. SPT order generalizes the notion of topological insulator to interacting systems.
Some suggest that topological order (or more precisely, string-net condensation) in local bosonic (spin) models has the potential to provide a unified origin for photons, electrons and other elementary particles in our universe.
| Physical sciences | Quantum mechanics | Physics |
3087674 | https://en.wikipedia.org/wiki/Plymouth%20Rock%20chicken | Plymouth Rock chicken | The Plymouth Rock is an American breed of domestic chicken. It was first seen in Massachusetts in the nineteenth century and for much of the early twentieth century was the most widely kept chicken breed in the United States. It is a dual-purpose bird, raised both for its meat and for its brown eggs. It is resistant to cold, easy to manage, and a good sitter.
History
The Plymouth Rock was first shown in Boston in 1849, but was then not seen for another twenty years. In 1869, in Worcester, Massachusetts, one D.A. Upham cross-bred some Black Java hens with a cock with barred plumage and a single comb; he selectively bred for barred plumage and clean (featherless) legs. His birds were shown in Worcester in 1869; the modern Plymouth Rock is thought to derive from them. Other people have been associated with the development of the Plymouth Rock, as have other chicken breeds including the Brahma, the Cochin (both white and buff), the Dominique and the White-faced Black Spanish. According to the Livestock Conservancy, it may have originated from cross-breeding of Java birds with single-combed Dominiques; or, based on genomic analysis, principally from the Dominique, with substantial contribution from the Java and Cochin and some input from other breeds.
The Plymouth Rock was included in the first edition of the American Standard of Perfection of the new American Poultry Association in 1874. The barred plumage pattern was the original one; other colors were later added.
It became the most widespread chicken breed in the United States and remained so until about the time of World War II. With the advent of industrial chicken farming, it was much used in the development of broiler hybrids but began to fall in popularity as a domestic fowl.
In 2023 the Plymouth Rock was listed by the Livestock Conservancy as 'recovering', meaning that there were at least new registrations per year. Worldwide, numbers for the Plymouth Rock are reported at almost ; about are reported for the Barred Plymouth Rock and over for the White variety.
Characteristics
The Plymouth Rock is easy to manage, is early-feathering, has good resistance to cold and is a good sitter. It has a single comb with five points; the comb, wattles and ear-lobes are bright red. The legs are yellow and unfeathered. The beak is yellow or horn-colored. The back is long and broad, and the breast fairly deep.
In the United States, seven color varieties of the Plymouth Rock are recognized: barred, blue, buff, Columbian, partridge, silver-penciled and white. Ten plumage varieties are listed by the Entente Européenne d'Aviculture et de Cuniculture, of which five – the barred, black, buff, Columbian and white – are recognized by the Poultry Club of Great Britain. In Australia, the barred variant is split into two separate colors, dark barred and light barred.
Use
The Plymouth Rock is a dual-purpose breed and is kept both for its meat and for its large brown eggs, of which it lays about 200 per year. The eggs weigh about .
In industrial agriculture, crosses of suitable strains of white Plymouth Rock with industrial strains of white Cornish constitute the principal stock of American broiler production.
| Biology and health sciences | Chickens | Animals |
17481271 | https://en.wikipedia.org/wiki/Fluorine | Fluorine | Fluorine is a chemical element; it has symbol F and atomic number 9. It is the lightest halogen and exists at standard conditions as pale yellow diatomic gas. Fluorine is extremely reactive as it reacts with all other elements except for the light inert gases. It is highly toxic.
Among the elements, fluorine ranks 24th in cosmic abundance and 13th in crustal abundance. Fluorite, the primary mineral source of fluorine, which gave the element its name, was first described in 1529; as it was added to metal ores to lower their melting points for smelting, the Latin verb meaning gave the mineral its name. Proposed as an element in 1810, fluorine proved difficult and dangerous to separate from its compounds, and several early experimenters died or sustained injuries from their attempts. Only in 1886 did French chemist Henri Moissan isolate elemental fluorine using low-temperature electrolysis, a process still employed for modern production. Industrial production of fluorine gas for uranium enrichment, its largest application, began during the Manhattan Project in World War II.
Owing to the expense of refining pure fluorine, most commercial applications use fluorine compounds, with about half of mined fluorite used in steelmaking. The rest of the fluorite is converted into hydrogen fluoride en route to various organic fluorides, or into cryolite, which plays a key role in aluminium refining. The carbon–fluorine bond is usually very stable. Organofluorine compounds are widely used as refrigerants, electrical insulation, and PTFE (Teflon). Pharmaceuticals such as atorvastatin and fluoxetine contain C−F bonds. The fluoride ion from dissolved fluoride salts inhibits dental cavities and so finds use in toothpaste and water fluoridation. Global fluorochemical sales amount to more than US$15 billion a year.
Fluorocarbon gases are generally greenhouse gases with global-warming potentials 100 to 23,500 times that of carbon dioxide, and SF6 has the highest global warming potential of any known substance. Organofluorine compounds often persist in the environment due to the strength of the carbon–fluorine bond. Fluorine has no known metabolic role in mammals; a few plants and marine sponges synthesize organofluorine poisons (most often monofluoroacetates) that help deter predation.
Characteristics
Electron configuration
Fluorine atoms have nine electrons, one fewer than neon, and electron configuration 1s22s22p5: two electrons in a filled inner shell and seven in an outer shell requiring one more to be filled. The outer electrons are ineffective at nuclear shielding, and experience a high effective nuclear charge of 9 − 2 = 7; this affects the atom's physical properties.
Fluorine's first ionization energy is third-highest among all elements, behind helium and neon, which complicates the removal of electrons from neutral fluorine atoms. It also has a high electron affinity, second only to chlorine, and tends to capture an electron to become isoelectronic with the noble gas neon; it has the highest electronegativity of any reactive element. Fluorine atoms have a small covalent radius of around 60 picometers, similar to those of its period neighbors oxygen and neon.
Reactivity
The bond energy of difluorine is much lower than that of either or and similar to the easily cleaved peroxide bond; this, along with high electronegativity, accounts for fluorine's easy dissociation, high reactivity, and strong bonds to non-fluorine atoms. Conversely, bonds to other atoms are very strong because of fluorine's high electronegativity. Unreactive substances like powdered steel, glass fragments, and asbestos fibers react quickly with cold fluorine gas; wood and water spontaneously combust under a fluorine jet.
Reactions of elemental fluorine with metals require varying conditions. Alkali metals cause explosions and alkaline earth metals display vigorous activity in bulk; to prevent passivation from the formation of metal fluoride layers, most other metals such as aluminium and iron must be powdered, and noble metals require pure fluorine gas at . Some solid nonmetals (sulfur, phosphorus) react vigorously in liquid fluorine. Hydrogen sulfide and sulfur dioxide combine readily with fluorine, the latter sometimes explosively; sulfuric acid exhibits much less activity, requiring elevated temperatures.
Hydrogen, like some of the alkali metals, reacts explosively with fluorine. Carbon, as lamp black, reacts at room temperature to yield tetrafluoromethane. Graphite combines with fluorine above to produce non-stoichiometric carbon monofluoride; higher temperatures generate gaseous fluorocarbons, sometimes with explosions. Carbon dioxide and carbon monoxide react at or just above room temperature, whereas paraffins and other organic chemicals generate strong reactions: even completely substituted haloalkanes such as carbon tetrachloride, normally incombustible, may explode. Although nitrogen trifluoride is stable, nitrogen requires an electric discharge at elevated temperatures for reaction with fluorine to occur, due to the very strong triple bond in elemental nitrogen; ammonia may react explosively. Oxygen does not combine with fluorine under ambient conditions, but can be made to react using electric discharge at low temperatures and pressures; the products tend to disintegrate into their constituent elements when heated. Heavier halogens react readily with fluorine as does the noble gas radon; of the other noble gases, only xenon and krypton react, and only under special conditions. Argon does not react with fluorine gas; however, it does form a compound with fluorine, argon fluorohydride.
Phases
At room temperature, fluorine is a gas of diatomic molecules, pale yellow when pure (sometimes described as yellow-green). It has a characteristic halogen-like pungent and biting odor detectable at 20 ppb. Fluorine condenses into a bright yellow liquid at , a transition temperature similar to those of oxygen and nitrogen.
Fluorine has two solid forms, α- and β-fluorine. The latter crystallizes at and is transparent and soft, with the same disordered cubic structure of freshly crystallized solid oxygen, unlike the orthorhombic systems of other solid halogens. Further cooling to induces a phase transition into opaque and hard α-fluorine, which has a monoclinic structure with dense, angled layers of molecules. The transition from β- to α-fluorine is more exothermic than the condensation of fluorine, and can be violent.
Isotopes
Only one isotope of fluorine occurs naturally in abundance, the stable isotope . It has a high magnetogyric ratio and exceptional sensitivity to magnetic fields; because it is also the only stable isotope, it is used in magnetic resonance imaging. Eighteen radioisotopes with mass numbers 13–31 have been synthesized, of which is the most stable with a half-life of 109.734 minutes. is a natural trace radioisotope produced by cosmic ray spallation of atmospheric argon as well as by reaction of protons with natural oxygen: 18O + p → 18F + n. Other radioisotopes have half-lives less than 70 seconds; most decay in less than half a second. The isotopes and undergo β+ decay and electron capture, lighter isotopes decay by proton emission, and those heavier than undergo β− decay (the heaviest ones with delayed neutron emission). Two metastable isomers of fluorine are known, , with a half-life of 162(7) nanoseconds, and , with a half-life of 2.2(1) milliseconds.
Occurrence
Universe
Among the lighter elements, fluorine's abundance value of 400 ppb (parts per billion) – 24th among elements in the universe – is exceptionally low: other elements from carbon to magnesium are twenty or more times as common. This is because stellar nucleosynthesis processes bypass fluorine, and any fluorine atoms otherwise created have high nuclear cross sections, allowing collisions with hydrogen or helium to generate oxygen or neon respectively.
Beyond this transient existence, three explanations have been proposed for the presence of fluorine:
during type II supernovae, bombardment of neon atoms by neutrinos could transmute them to fluorine;
the solar wind of Wolf–Rayet stars could blow fluorine away from any hydrogen or helium atoms; or
fluorine is borne out on convection currents arising from fusion in asymptotic giant branch stars.
Earth
Fluorine is the 13th most abundant element in Earth's crust at 600–700 ppm (parts per million) by mass. Though believed not to occur naturally, elemental fluorine has been shown to be present as an occlusion in antozonite, a variant of fluorite. Most fluorine exists as fluoride-containing minerals. Fluorite, fluorapatite and cryolite are the most industrially significant. Fluorite (), also known as fluorspar, abundant worldwide, is the main source of fluoride, and hence fluorine. China and Mexico are the major suppliers. Fluorapatite (Ca5(PO4)3F), which contains most of the world's fluoride, is an inadvertent source of fluoride as a byproduct of fertilizer production. Cryolite (), used in the production of aluminium, is the most fluorine-rich mineral. Economically viable natural sources of cryolite have been exhausted, and most is now synthesised commercially.
Other minerals such as topaz contain fluorine. Fluorides, unlike other halides, are insoluble and do not occur in commercially favorable concentrations in saline waters. Trace quantities of organofluorines of uncertain origin have been detected in volcanic eruptions and geothermal springs. The existence of gaseous fluorine in crystals, suggested by the smell of crushed antozonite, is contentious; a 2012 study reported the presence of 0.04% by weight in antozonite, attributing these inclusions to radiation from the presence of tiny amounts of uranium.
History
Early discoveries
In 1529, Georgius Agricola described fluorite as an additive used to lower the melting point of metals during smelting. He penned the Latin word fluorēs (fluor, flow) for fluorite rocks. The name later evolved into fluorspar (still commonly used) and then fluorite. The composition of fluorite was later determined to be calcium difluoride.
Hydrofluoric acid was used in glass etching from 1720 onward. Andreas Sigismund Marggraf first characterized it in 1764 when he heated fluorite with sulfuric acid, and the resulting solution corroded its glass container. Swedish chemist Carl Wilhelm Scheele repeated the experiment in 1771, and named the acidic product fluss-spats-syran (fluorspar acid). In 1810, the French physicist André-Marie Ampère suggested that hydrogen and an element analogous to chlorine constituted hydrofluoric acid. He also proposed in a letter to Sir Humphry Davy dated August 26, 1812 that this then-unknown substance may be named fluorine from fluoric acid and the -ine suffix of other halogens. This word, often with modifications, is used in most European languages; however, Greek, Russian, and some others, following Ampère's later suggestion, use the name ftor or derivatives, from the Greek φθόριος (phthorios, destructive). The New Latin name fluorum gave the element its current symbol F; Fl was used in early papers.
Isolation
Initial studies on fluorine were so dangerous that several 19th-century experimenters were deemed "fluorine martyrs" after misfortunes with hydrofluoric acid. Isolation of elemental fluorine was hindered by the extreme corrosiveness of both elemental fluorine itself and hydrogen fluoride, as well as the lack of a simple and suitable electrolyte. Edmond Frémy postulated that electrolysis of pure hydrogen fluoride to generate fluorine was feasible and devised a method to produce anhydrous samples from acidified potassium bifluoride; instead, he discovered that the resulting (dry) hydrogen fluoride did not conduct electricity. Frémy's former student Henri Moissan persevered, and after much trial and error found that a mixture of potassium bifluoride and dry hydrogen fluoride was a conductor, enabling electrolysis. To prevent rapid corrosion of the platinum in his electrochemical cells, he cooled the reaction to extremely low temperatures in a special bath and forged cells from a more resistant mixture of platinum and iridium, and used fluorite stoppers. In 1886, after 74 years of effort by many chemists, Moissan isolated elemental fluorine.
In 1906, two months before his death, Moissan received the Nobel Prize in Chemistry, with the following citation:
Later uses
The Frigidaire division of General Motors (GM) experimented with chlorofluorocarbon refrigerants in the late 1920s, and Kinetic Chemicals was formed as a joint venture between GM and DuPont in 1930 hoping to market Freon-12 () as one such refrigerant. It replaced earlier and more toxic compounds, increased demand for kitchen refrigerators, and became profitable; by 1949 DuPont had bought out Kinetic and marketed several other Freon compounds. Polytetrafluoroethylene (Teflon) was serendipitously discovered in 1938 by Roy J. Plunkett while working on refrigerants at Kinetic, and its superlative chemical and thermal resistance lent it to accelerated commercialization and mass production by 1941.
Large-scale production of elemental fluorine began during World War II. Germany used high-temperature electrolysis to make tons of the planned incendiary chlorine trifluoride and the Manhattan Project used huge quantities to produce uranium hexafluoride for uranium enrichment. Since is as corrosive as fluorine, gaseous diffusion plants required special materials: nickel for membranes, fluoropolymers for seals, and liquid fluorocarbons as coolants and lubricants. This burgeoning nuclear industry later drove post-war fluorochemical development.
Compounds
Fluorine has a rich chemistry, encompassing organic and inorganic domains. It combines with metals, nonmetals, metalloids, and most noble gases. Fluorine's high electron affinity results in a preference for ionic bonding; when it forms covalent bonds, these are polar, and almost always single.
Oxidation states
In compounds, fluorine almost exclusively assumes an oxidation state of −1. Fluorine in is defined to have oxidation state 0. The unstable species and , which decompose at around 40 K, have intermediate oxidation states; and a few related species are predicted to be stable.
Metals
Alkali metals form ionic and highly soluble monofluorides; these have the cubic arrangement of sodium chloride and analogous chlorides. Alkaline earth difluorides possess strong ionic bonds but are insoluble in water, with the exception of beryllium difluoride, which also exhibits some covalent character and has a quartz-like structure. Rare earth elements and many other metals form mostly ionic trifluorides.
Covalent bonding first comes to prominence in the tetrafluorides: those of zirconium, hafnium and several actinides are ionic with high melting points, while those of titanium, vanadium, and niobium are polymeric, melting or decomposing at no more than . Pentafluorides continue this trend with their linear polymers and oligomeric complexes. Thirteen metal hexafluorides are known, all octahedral, and are mostly volatile solids but for liquid and , and gaseous . Rhenium heptafluoride, the only characterized metal heptafluoride, is a low-melting molecular solid with pentagonal bipyramidal molecular geometry. Metal fluorides with more fluorine atoms are particularly reactive.
Hydrogen
Hydrogen and fluorine combine to yield hydrogen fluoride, in which discrete molecules form clusters by hydrogen bonding, resembling water more than hydrogen chloride. It boils at a much higher temperature than heavier hydrogen halides and unlike them is miscible with water. Hydrogen fluoride readily hydrates on contact with water to form aqueous hydrogen fluoride, also known as hydrofluoric acid. Unlike the other hydrohalic acids, which are strong, hydrofluoric acid is a weak acid at low concentrations. However, it can attack glass, something the other acids cannot do.
Other reactive nonmetals
Binary fluorides of metalloids and p-block nonmetals are generally covalent and volatile, with varying reactivities. Period 3 and heavier nonmetals can form hypervalent fluorides.
Boron trifluoride is planar and possesses an incomplete octet. It functions as a Lewis acid and combines with Lewis bases like ammonia to form adducts. Carbon tetrafluoride is tetrahedral and inert; its group analogues, silicon and germanium tetrafluoride, are also tetrahedral but behave as Lewis acids. The pnictogens form trifluorides that increase in reactivity and basicity with higher molecular weight, although nitrogen trifluoride resists hydrolysis and is not basic. The pentafluorides of phosphorus, arsenic, and antimony are more reactive than their respective trifluorides, with antimony pentafluoride the strongest neutral Lewis acid known, only behind gold pentafluoride.
Chalcogens have diverse fluorides: unstable difluorides have been reported for oxygen (the only known compound with oxygen in an oxidation state of +2), sulfur, and selenium; tetrafluorides and hexafluorides exist for sulfur, selenium, and tellurium. The latter are stabilized by more fluorine atoms and lighter central atoms, so sulfur hexafluoride is especially inert. Chlorine, bromine, and iodine can each form mono-, tri-, and pentafluorides, but only iodine heptafluoride has been characterized among possible interhalogen heptafluorides. Many of them are powerful sources of fluorine atoms, and industrial applications using chlorine trifluoride require precautions similar to those using fluorine.
Noble gases
Noble gases, having complete electron shells, defied reaction with other elements until 1962 when Neil Bartlett reported synthesis of xenon hexafluoroplatinate; xenon difluoride, tetrafluoride, hexafluoride, and multiple oxyfluorides have been isolated since then. Among other noble gases, krypton forms a difluoride, and radon and fluorine generate a solid suspected to be radon difluoride. Binary fluorides of lighter noble gases are exceptionally unstable: argon and hydrogen fluoride combine under extreme conditions to give argon fluorohydride. Helium has no long-lived fluorides, and no neon fluoride has ever been observed; helium fluorohydride has been detected for milliseconds at high pressures and low temperatures.
Organic compounds
The carbon–fluorine bond is organic chemistry's strongest, and gives stability to organofluorines. It is almost non-existent in nature, but is used in artificial compounds. Research in this area is usually driven by commercial applications; the compounds involved are diverse and reflect the complexity inherent in organic chemistry.
Discrete molecules
The substitution of hydrogen atoms in an alkane by progressively more fluorine atoms gradually alters several properties: melting and boiling points are lowered, density increases, solubility in hydrocarbons decreases and overall stability increases. Perfluorocarbons, in which all hydrogen atoms are substituted, are insoluble in most organic solvents, reacting at ambient conditions only with sodium in liquid ammonia.
The term perfluorinated compound is used for what would otherwise be a perfluorocarbon if not for the presence of a functional group, often a carboxylic acid. These compounds share many properties with perfluorocarbons such as stability and hydrophobicity, while the functional group augments their reactivity, enabling them to adhere to surfaces or act as surfactants. Fluorosurfactants, in particular, can lower the surface tension of water more than their hydrocarbon-based analogues. Fluorotelomers, which have some unfluorinated carbon atoms near the functional group, are also regarded as perfluorinated.
Polymers
Polymers exhibit the same stability increases afforded by fluorine substitution (for hydrogen) in discrete molecules; their melting points generally increase too. Polytetrafluoroethylene (PTFE), the simplest fluoropolymer and perfluoro analogue of polyethylene with structural unit ––, demonstrates this change as expected, but its very high melting point makes it difficult to mold. Various PTFE derivatives are less temperature-tolerant but easier to mold: fluorinated ethylene propylene replaces some fluorine atoms with trifluoromethyl groups, perfluoroalkoxy alkanes do the same with trifluoromethoxy groups, and Nafion contains perfluoroether side chains capped with sulfonic acid groups. Other fluoropolymers retain some hydrogen atoms; polyvinylidene fluoride has half the fluorine atoms of PTFE and polyvinyl fluoride has a quarter, but both behave much like perfluorinated polymers.
Production
Elemental fluorine and virtually all fluorine compounds are produced from hydrogen fluoride or its aqueous solution, hydrofluoric acid. Hydrogen fluoride is produced in kilns by the endothermic reaction of fluorite (CaF2) with sulfuric acid:
CaF2 + H2SO4 → 2 HF(g) + CaSO4
The gaseous HF can then be absorbed in water or liquefied.
About 20% of manufactured HF is a byproduct of fertilizer production, which produces hexafluorosilicic acid (H2SiF6), which can be degraded to release HF thermally and by hydrolysis:
H2SiF6 → 2 HF + SiF4
SiF4 + 2 H2O → 4 HF + SiO2
Industrial routes to F2
Moissan's method is used to produce industrial quantities of fluorine, via the electrolysis of a potassium bifluoride/hydrogen fluoride mixture: hydrogen ions are reduced at a steel container cathode and fluoride ions are oxidized at a carbon block anode, under 8–12 volts, to generate hydrogen and fluorine gas respectively. Temperatures are elevated, KF•2HF melting at and being electrolyzed at . KF, which acts to provide electrical conductivity, is essential since pure HF cannot be electrolyzed because it is virtually non-conductive. Fluorine can be stored in steel cylinders that have passivated interiors, at temperatures below ; otherwise nickel can be used. Regulator valves and pipework are made of nickel, the latter possibly using Monel instead. Frequent passivation, along with the strict exclusion of water and greases, must be undertaken. In the laboratory, glassware may carry fluorine gas under low pressure and anhydrous conditions; some sources instead recommend nickel-Monel-PTFE systems.
Laboratory routes
While preparing for a 1986 conference to celebrate the centennial of Moissan's achievement, Karl O. Christe reasoned that chemical fluorine generation should be feasible since some metal fluoride anions have no stable neutral counterparts; their acidification potentially triggers oxidation instead. He devised a method which evolves fluorine at high yield and atmospheric pressure:
2 KMnO4 + 2 KF + 10 HF + 3 H2O2 → 2 K2MnF6 + 8 H2O + 3 O2↑
2 K2MnF6 + 4 SbF5 → 4 KSbF6 + 2 MnF3 + F2↑
Christe later commented that the reactants "had been known for more than 100 years and even Moissan could have come up with this scheme." As late as 2008, some references still asserted that fluorine was too reactive for any chemical isolation.
Industrial applications
Fluorite mining, which supplies most global fluorine, peaked in 1989 when 5.6 million metric tons of ore were extracted. Chlorofluorocarbon restrictions lowered this to 3.6 million tons in 1994; production has since been increasing. Around 4.5 million tons of ore and revenue of US$550 million were generated in 2003; later reports estimated 2011 global fluorochemical sales at $15 billion and predicted 2016–18 production figures of 3.5 to 5.9 million tons, and revenue of at least $20 billion. Froth flotation separates mined fluorite into two main metallurgical grades of equal proportion: 60–85% pure metspar is almost all used in iron smelting whereas 97%+ pure acidspar is mainly converted to the key industrial intermediate hydrogen fluoride.
At least 17,000 metric tons of fluorine are produced each year. It costs only $5–8 per kilogram as uranium or sulfur hexafluoride, but many times more as an element because of handling challenges. Most processes using free fluorine in large amounts employ in situ generation under vertical integration.
The largest application of fluorine gas, consuming up to 7,000 metric tons annually, is in the preparation of for the nuclear fuel cycle. Fluorine is used to fluorinate uranium tetrafluoride, itself formed from uranium dioxide and hydrofluoric acid. Fluorine is monoisotopic, so any mass differences between molecules are due to the presence of or , enabling uranium enrichment via gaseous diffusion or gas centrifuge. About 6,000 metric tons per year go into producing the inert dielectric for high-voltage transformers and circuit breakers, eliminating the need for hazardous polychlorinated biphenyls associated with devices. Several fluorine compounds are used in electronics: rhenium and tungsten hexafluoride in chemical vapor deposition, tetrafluoromethane in plasma etching and nitrogen trifluoride in cleaning equipment. Fluorine is also used in the synthesis of organic fluorides, but its reactivity often necessitates conversion first to the gentler , , or , which together allow calibrated fluorination. Fluorinated pharmaceuticals use sulfur tetrafluoride instead.
Inorganic fluorides
As with other iron alloys, around metspar is added to each metric ton of steel; the fluoride ions lower its melting point and viscosity. Alongside its role as an additive in materials like enamels and welding rod coats, most acidspar is reacted with sulfuric acid to form hydrofluoric acid, which is used in steel pickling, glass etching and alkane cracking. One-third of HF goes into synthesizing cryolite and aluminium trifluoride, both fluxes in the Hall–Héroult process for aluminium extraction; replenishment is necessitated by their occasional reactions with the smelting apparatus. Each metric ton of aluminium requires about of flux. Fluorosilicates consume the second largest portion, with sodium fluorosilicate used in water fluoridation and laundry effluent treatment, and as an intermediate en route to cryolite and silicon tetrafluoride. Other important inorganic fluorides include those of cobalt, nickel, and ammonium.
Organic fluorides
Organofluorides consume over 20% of mined fluorite and over 40% of hydrofluoric acid, with refrigerant gases dominating and fluoropolymers increasing their market share. Surfactants are a minor application but generate over $1 billion in annual revenue. Due to the danger from direct hydrocarbon–fluorine reactions above , industrial fluorocarbon production is indirect, mostly through halogen exchange reactions such as Swarts fluorination, in which chlorocarbon chlorines are substituted for fluorines by hydrogen fluoride under catalysts. Electrochemical fluorination subjects hydrocarbons to electrolysis in hydrogen fluoride, and the Fowler process treats them with solid fluorine carriers like cobalt trifluoride.
Refrigerant gases
Halogenated refrigerants, termed Freons in informal contexts, are identified by R-numbers that denote the amount of fluorine, chlorine, carbon, and hydrogen present. Chlorofluorocarbons (CFCs) like R-11, R-12, and R-114 once dominated organofluorines, peaking in production in the 1980s. Used for air conditioning systems, propellants and solvents, their production was below one-tenth of this peak by the early 2000s, after widespread international prohibition. Hydrochlorofluorocarbons (HCFCs) and hydrofluorocarbons (HFCs) were designed as replacements; their synthesis consumes more than 90% of the fluorine in the organic industry. Important HCFCs include R-22, chlorodifluoromethane, and R-141b. The main HFC is R-134a with a new type of molecule HFO-1234yf, a Hydrofluoroolefin (HFO) coming to prominence owing to its global warming potential of less than 1% that of HFC-134a.
Polymers
About 180,000 metric tons of fluoropolymers were produced in 2006 and 2007, generating over $3.5 billion revenue per year. The global market was estimated at just under $6 billion in 2011. Fluoropolymers can only be formed by polymerizing free radicals.
Polytetrafluoroethylene (PTFE), sometimes called by its DuPont name Teflon, represents 60–80% by mass of the world's fluoropolymer production. The largest application is in electrical insulation since PTFE is an excellent dielectric. It is also used in the chemical industry where corrosion resistance is needed, in coating pipes, tubing, and gaskets. Another major use is in PFTE-coated fiberglass cloth for stadium roofs. The major consumer application is for non-stick cookware. Jerked PTFE film becomes expanded PTFE (ePTFE), a fine-pored membrane sometimes referred to by the brand name Gore-Tex and used for rainwear, protective apparel, and filters; ePTFE fibers may be made into seals and dust filters. Other fluoropolymers, including fluorinated ethylene propylene, mimic PTFE's properties and can substitute for it; they are more moldable, but also more costly and have lower thermal stability. Films from two different fluoropolymers replace glass in solar cells.
The chemically resistant (but expensive) fluorinated ionomers are used as electrochemical cell membranes, of which the first and most prominent example is Nafion. Developed in the 1960s, it was initially deployed as fuel cell material in spacecraft and then replaced mercury-based chloralkali process cells. Recently, the fuel cell application has reemerged with efforts to install proton exchange membrane fuel cells into automobiles. Fluoroelastomers such as Viton are crosslinked fluoropolymer mixtures mainly used in O-rings; perfluorobutane (C4F10) is used as a fire-extinguishing agent.
Surfactants
Fluorosurfactants are small organofluorine molecules used for repelling water and stains. Although expensive (comparable to pharmaceuticals at $200–2000 per kilogram), they yielded over $1 billion in annual revenues by 2006; Scotchgard alone generated over $300 million in 2000. Fluorosurfactants are a minority in the overall surfactant market, most of which is taken up by much cheaper hydrocarbon-based products. Applications in paints are burdened by compounding costs; this use was valued at only $100 million in 2006.
Agrichemicals
About 30% of agrichemicals contain fluorine, most of them herbicides and fungicides with a few crop regulators. Fluorine substitution, usually of a single atom or at most a trifluoromethyl group, is a robust modification with effects analogous to fluorinated pharmaceuticals: increased biological stay time, membrane crossing, and altering of molecular recognition. Trifluralin is a prominent example, with large-scale use in the U.S. as a weedkiller, but it is a suspected carcinogen and has been banned in many European countries. Sodium monofluoroacetate (1080) is a mammalian poison in which one sodium acetate hydrogen is replaced with fluorine; it disrupts cell metabolism by replacing acetate in the citric acid cycle. First synthesized in the late 19th century, it was recognized as an insecticide in the early 20th century, and was later deployed in its current use. New Zealand, the largest consumer of 1080, uses it to protect kiwis from the invasive Australian common brushtail possum. Europe and the U.S. have banned 1080.
Medicinal applications
Dental care
Population studies from the mid-20th century onwards show topical fluoride reduces dental caries. This was first attributed to the conversion of tooth enamel hydroxyapatite into the more durable fluorapatite, but studies on pre-fluoridated teeth refuted this hypothesis, and current theories involve fluoride aiding enamel growth in small caries. After studies of children in areas where fluoride was naturally present in drinking water, controlled public water supply fluoridation to fight tooth decay began in the 1940s and is now applied to water supplying 6 percent of the global population, including two-thirds of Americans. Reviews of the scholarly literature in 2000 and 2007 associated water fluoridation with a significant reduction of tooth decay in children. Despite such endorsements and evidence of no adverse effects other than mostly benign dental fluorosis, opposition still exists on ethical and safety grounds. The benefits of fluoridation have lessened, possibly due to other fluoride sources, but are still measurable in low-income groups. Sodium monofluorophosphate and sometimes sodium or tin(II) fluoride are often found in fluoride toothpastes, first introduced in the U.S. in 1955 and now ubiquitous in developed countries, alongside fluoridated mouthwashes, gels, foams, and varnishes.
Pharmaceuticals
Twenty percent of modern pharmaceuticals contain fluorine. One of these, the cholesterol-reducer atorvastatin (Lipitor), made more revenue than any other drug until it became generic in 2011. The combination asthma prescription Seretide, a top-ten revenue drug in the mid-2000s, contains two active ingredients, one of which – fluticasone – is fluorinated. Many drugs are fluorinated to delay inactivation and lengthen dosage periods because the carbon–fluorine bond is very stable. Fluorination also increases lipophilicity because the bond is more hydrophobic than the carbon–hydrogen bond, and this often helps in cell membrane penetration and hence bioavailability.
Tricyclics and other pre-1980s antidepressants had several side effects due to their non-selective interference with neurotransmitters other than the serotonin target; the fluorinated fluoxetine was selective and one of the first to avoid this problem. Many current antidepressants receive this same treatment, including the selective serotonin reuptake inhibitors: citalopram, its enantiomer escitalopram, and fluvoxamine and paroxetine. Quinolones are artificial broad-spectrum antibiotics that are often fluorinated to enhance their effects. These include ciprofloxacin and levofloxacin. Fluorine also finds use in steroids: fludrocortisone is a blood pressure-raising mineralocorticoid, and triamcinolone and dexamethasone are strong glucocorticoids. The majority of inhaled anesthetics are heavily fluorinated; the prototype halothane is much more inert and potent than its contemporaries. Later compounds such as the fluorinated ethers sevoflurane and desflurane are better than halothane and are almost insoluble in blood, allowing faster waking times.
PET scanning
Fluorine-18 is often found in radioactive tracers for positron emission tomography, as its half-life of almost two hours is long enough to allow for its transport from production facilities to imaging centers. The most common tracer is fluorodeoxyglucose which, after intravenous injection, is taken up by glucose-requiring tissues such as the brain and most malignant tumors; computer-assisted tomography can then be used for detailed imaging.
Oxygen carriers
Liquid fluorocarbons can hold large volumes of oxygen or carbon dioxide, more so than blood, and have attracted attention for their possible uses in artificial blood and in liquid breathing. Because fluorocarbons do not normally mix with water, they must be mixed into emulsions (small droplets of perfluorocarbon suspended in water) to be used as blood. One such product, Oxycyte, has been through initial clinical trials. These substances can aid endurance athletes and are banned from sports; one cyclist's near death in 1998 prompted an investigation into their abuse. Applications of pure perfluorocarbon liquid breathing (which uses pure perfluorocarbon liquid, not a water emulsion) include assisting burn victims and premature babies with deficient lungs. Partial and complete lung filling have been considered, though only the former has had any significant tests in humans. An Alliance Pharmaceuticals effort reached clinical trials but was abandoned because the results were not better than normal therapies.
Biological role
Fluorine is not essential for humans and other mammals, but small amounts are known to be beneficial for the strengthening of dental enamel (where the formation of fluorapatite makes the enamel more resistant to attack, from acids produced by bacterial fermentation of sugars). Small amounts of fluorine may be beneficial for bone strength, but the latter has not been definitively established. Both the WHO and the Institute of Medicine of the US National Academies publish recommended daily allowance (RDA) and upper tolerated intake of fluorine, which varies with age and gender.
Natural organofluorines have been found in microorganisms, plants and, recently, animals. The most common is fluoroacetate, which is used as a defense against herbivores by at least 40 plants in Africa, Australia and Brazil. Other examples include terminally fluorinated fatty acids, fluoroacetone, and 2-fluorocitrate. An enzyme that binds fluorine to carbon – adenosyl-fluoride synthase – was discovered in bacteria in 2002.
Toxicity
Elemental fluorine is highly toxic to living organisms. Its effects in humans start at concentrations lower than hydrogen cyanide's 50 ppm and are similar to those of chlorine: significant irritation of the eyes and respiratory system as well as liver and kidney damage occur above 25 ppm, which is the immediately dangerous to life and health value for fluorine. The eyes and nose are seriously damaged at 100 ppm, and inhalation of 1,000 ppm fluorine will cause death in minutes, compared to 270 ppm for hydrogen cyanide.
Hydrofluoric acid
Hydrofluoric acid is the weakest of the hydrohalic acids, having a pKa of 3.2 at 25 °C. Pure hydrogen fluoride is a volatile liquid due to the presence of hydrogen bonding, while the other hydrogen halides are gases. It is able to attack glass, concrete, metals, and organic matter.
Hydrofluoric acid is a contact poison with greater hazards than many strong acids like sulfuric acid even though it is weak: it remains neutral in aqueous solution and thus penetrates tissue faster, whether through inhalation, ingestion or the skin, and at least nine U.S. workers died in such accidents from 1984 to 1994. It reacts with calcium and magnesium in the blood leading to hypocalcemia and possible death through cardiac arrhythmia. Insoluble calcium fluoride formation triggers strong pain and burns larger than 160 cm2 (25 in2) can cause serious systemic toxicity.
Exposure may not be evident for eight hours for 50% HF, rising to 24 hours for lower concentrations, and a burn may initially be painless as hydrogen fluoride affects nerve function. If skin has been exposed to HF, damage can be reduced by rinsing it under a jet of water for 10–15 minutes and removing contaminated clothing. Calcium gluconate is often applied next, providing calcium ions to bind with fluoride; skin burns can be treated with 2.5% calcium gluconate gel or special rinsing solutions. Hydrofluoric acid absorption requires further medical treatment; calcium gluconate may be injected or administered intravenously. Using calcium chloride – a common laboratory reagent – in lieu of calcium gluconate is contraindicated, and may lead to severe complications. Excision or amputation of affected parts may be required.
Fluoride ion
Soluble fluorides are moderately toxic: 5–10 g sodium fluoride, or 32–64 mg fluoride ions per kilogram of body mass, represents a lethal dose for adults. One-fifth of the lethal dose can cause adverse health effects, and chronic excess consumption may lead to skeletal fluorosis, which affects millions in Asia and Africa, and, in children, to reduced intelligence. Ingested fluoride forms hydrofluoric acid in the stomach which is easily absorbed by the intestines, where it crosses cell membranes, binds with calcium and interferes with various enzymes, before urinary excretion. Exposure limits are determined by urine testing of the body's ability to clear fluoride ions.
Historically, most cases of fluoride poisoning have been caused by accidental ingestion of insecticides containing inorganic fluorides. Most current calls to poison control centers for possible fluoride poisoning come from the ingestion of fluoride-containing toothpaste. Malfunctioning water fluoridation equipment is another cause: one incident in Alaska affected almost 300 people and killed one person. Dangers from toothpaste are aggravated for small children, and the Centers for Disease Control and Prevention recommends supervising children below six brushing their teeth so that they do not swallow toothpaste. One regional study examined a year of pre-teen fluoride poisoning reports totaling 87 cases, including one death from ingesting insecticide. Most had no symptoms, but about 30% had stomach pains. A larger study across the U.S. had similar findings: 80% of cases involved children under six, and there were few serious cases.
Environmental concerns
Atmosphere
The Montreal Protocol, signed in 1987, set strict regulations on chlorofluorocarbons (CFCs) and bromofluorocarbons due to their ozone damaging potential (ODP). The high stability which suited them to their original applications also meant that they were not decomposing until they reached higher altitudes, where liberated chlorine and bromine atoms attacked ozone molecules. Even with the ban, and early indications of its efficacy, predictions warned that several generations would pass before full recovery. With one-tenth the ODP of CFCs, hydrochlorofluorocarbons (HCFCs) are the current replacements, and are themselves scheduled for substitution by 2030–2040 by hydrofluorocarbons (HFCs) with no chlorine and zero ODP. In 2007 this date was brought forward to 2020 for developed countries; the Environmental Protection Agency had already prohibited one HCFC's production and capped those of two others in 2003. Fluorocarbon gases are generally greenhouse gases with global-warming potentials (GWPs) of about 100 to 10,000; sulfur hexafluoride has a value of around 20,000. An outlier is HFO-1234yf which is a new type of refrigerant called a Hydrofluoroolefin (HFO) and has attracted global demand due to its GWP of less than 1 compared to 1,430 for the current refrigerant standard HFC-134a.
Biopersistence
Organofluorines exhibit biopersistence due to the strength of the carbon–fluorine bond. Perfluoroalkyl acids (PFAAs), which are sparingly water-soluble owing to their acidic functional groups, are noted persistent organic pollutants; perfluorooctanesulfonic acid (PFOS) and perfluorooctanoic acid (PFOA) are most often researched. PFAAs have been found in trace quantities worldwide from polar bears to humans, with PFOS and PFOA known to reside in breast milk and the blood of newborn babies. A 2013 review showed a slight correlation between groundwater and soil PFAA levels and human activity; there was no clear pattern of one chemical dominating, and higher amounts of PFOS were correlated to higher amounts of PFOA. In the body, PFAAs bind to proteins such as serum albumin; they tend to concentrate within humans in the liver and blood before excretion through the kidneys. Dwell time in the body varies greatly by species, with half-lives of days in rodents, and years in humans. High doses of PFOS and PFOA cause cancer and death in newborn rodents but human studies have not established an effect at current exposure levels.
| Physical sciences | Chemical elements_2 | null |
16234880 | https://en.wikipedia.org/wiki/Tree%20frog | Tree frog | A tree frog (or treefrog) is any species of frog that spends a major portion of its lifespan in trees, known as an arboreal state. Several lineages of frogs among the Neobatrachia suborder have given rise to treefrogs, although they are not closely related to each other.
Millions of years of convergent evolution have resulted in very similar morphology even in species that are not very closely related.
Furthermore, tree frogs in seasonally arid environments have adapted an extra-epidermal layer of lipid and mucus as an evolutionary convergent response to accommodate the periodic dehydration stress.
Description
As the name implies, these frogs are typically found in trees or other high-growing vegetation. They do not normally descend to the ground, except to mate and spawn, though some build foam nests on leaves and rarely leave the trees at all as adults, and Eleutherodactylus has evolved direct development and therefore does not need water for a tadpole stage.
Tree frogs are usually tiny as their weight has to be carried by the branches and twigs in their habitats. While some reach 10 cm (4 in) or more, they are typically smaller and more slender than terrestrial frogs. Tree frogs typically have well-developed discs at the finger and toe tips, they rely on several attachment mechanisms that vary with circumstances, tree frogs require static and dynamic, adhesive and frictional, reversible and repeatable force generation; the fingers and toes themselves, as well as the limbs, tend to be rather small, resulting in a superior grasping ability. The genus Chiromantis of the Rhacophoridae is most extreme in this respect: it can oppose two fingers to the other two, resulting in a vise-like grip.
Family
Tree frogs are members of these families or genera:
Hylidae, or "true" treefrogs, occur in the temperate to tropical parts of Eurasia north of the Himalayas, Australia and the Americas.
Rhacophoridae, or shrub frogs, are the treefrogs of tropical regions around the Indian Ocean: Africa, South Asia and Southeast Asia east to Lydekker's line. A few also occur in East Asia.
Centrolenidae, or glass frogs, are potentially closely related to hylids; these translucent frogs are native to Central and South America.
Hyperoliidae, or reed frogs, are closely related to the burrowing Microhylidae; these small frogs are native to sub-Saharan Africa.
Boophis is a genus of highly arboreal frogs which evolved from the toxic terrestrial Mantellidae of Madagascar.
Gallery
| Biology and health sciences | Frogs and toads | Animals |
16248198 | https://en.wikipedia.org/wiki/Rex%20rabbit | Rex rabbit | The term rex rabbit (without capitalization) refers informally to one of at least eight breeds of domestic rabbit. One such breed is the Rex, which is recognized by the American Rabbit Breeders Association (ARBA) and by the British Rabbit Council (BRC). Other modern-day rex rabbit breeds are listed below. Care must be taken to distinguish the rex rabbit breeds from the three types of rex rabbit fur for which they are known.
The Rex rabbit breed that is recognized by ARBA is a medium-sized rabbit with a commercial, round body and an ideal weight range of . The Rex has a slightly broader head than other breeds of rabbit, proportionate upright ears, and proportionally smaller feet. As with most larger breeds, the female (or doe) has a dewlap (a large flap of skin under the chin).
History and origin
The Rex is a variety of rabbit that exhibits plush fur that is often described as having a velvety texture. The breed originated in France in 1919. Its origin was a litter of wild gray rabbits and has been developed over the years by fanciers and the fur industry. The Rex Rabbit was first shown publicly at the Paris International Rabbit Show in 1924 and has been recognized as a standard breed in parts of Europe since 1925. Eugène Kohler at the University of Strasbourg developed the breed further, giving rise to colored varieties and the chinchilla rex, which is one of the main breeds used in the rabbit fur farming industry.
The Rex was first imported to the United States in 1924 following the Paris International Rabbit Show by American rabbit pioneer John C. Fehr and his partner Alfred Zimmerman.
Genetics
Many genes contribute to the Standard Rex breed. The definition of the breed is maintained by ARBA and the British Rabbit Council (BRC). The definition is based strictly on phenotype. The gross external features used to identify the rabbit include weight, coloration, coat texture and length. Of these features, amongst fanciers and the fur industry, the coat properties are of chief concern. This breed has a low to moderate activity level and can jump as high as three feet.
Fur
The word rex (with a lowercase 'r') refers to the unique characteristics of an animal's "rexed" fur. It is the result of a specific genetic mutation that is now deliberately sought in the Rex rabbit breeds. In rex fur, the hair protrudes outwards from the body, instead of lying flat, and the guard hairs are shortened to the length of the undercoat, or a bit longer.
Because of the density of hairs, rex fur is often described as plush or velvety. Three varieties of rex fur now exist: standard rex fur, short curly rex fur (as in the Astrex), and long curly rex fur (as in the Opossum).
Varieties
ARBA currently recognizes 16 color varieties for the Rex rabbit breed, while the BRC currently recognizes 32.
There are clubs for fanciers of particular rabbits and the Standard Rex is no exception.
Modern development
Currently, Rex Rabbits are kept as pets, as show rabbits, and for fur and meat production. The Rex remains the number one breed used in fur production for garments and toys as their fur lacks protruding guard hairs which in other breeds which must be shorn and plucked after the tanning process to resemble other animal furs. The bulk of furs produced in the United States are more of a by-product to the primary purpose of meat production and are of lesser quality. Meat production favors harvest of young animals (70–84 days) which is not sufficient time for the rabbit to grow a high quality adult fur. The fur is best after 6 months and during the coldest part of the year. The overwhelming number of furs produced in the USA are used in the creation of felt because the fur quality of commercially raised rabbits is of too low a quality for the garment trade. See: Production of Rabbit Skins and Hair for Textiles
Rex rabbit breeds
Astrex
Canadian Plush Lop
Mini Rex
Opossum
Plush Lop (Miniature)
Plush Lop (Standard)
Rex [US]
Velveteen Lop
| Biology and health sciences | Rabbits | Animals |
4217326 | https://en.wikipedia.org/wiki/Intracluster%20medium | Intracluster medium | In astronomy, the intracluster medium (ICM) is the superheated plasma that permeates a galaxy cluster. The gas consists mainly of ionized hydrogen and helium and accounts for most of the baryonic material in galaxy clusters. The ICM is heated to temperatures on the order of 10 to 100 megakelvins, emitting strong X-ray radiation.
Composition
The ICM is composed primarily of ordinary baryons, mainly ionised hydrogen and helium. This plasma is enriched with heavier elements, including iron. The average amount of heavier elements relative to hydrogen, known as metallicity in astronomy, ranges from a third to a half of the value in the sun. Studying the chemical composition of the ICMs as a function of radius has shown that cores of the galaxy clusters are more metal-rich than at larger radii. In some clusters (e.g. the Centaurus cluster) the metallicity of the gas can rise to above that of the sun. Due to the gravitational field of clusters, metal-enriched gas ejected from supernova remains gravitationally bound to the cluster as part of the ICM. By looking at varying redshift, which corresponds to looking at different epochs of the evolution of the Universe, the ICM can provide a history record of element production in a galaxy.
Roughly 15% of a galaxy cluster's mass resides in the ICM. The stars and galaxies contribute only around 5% to the total mass. It is theorized that most of the mass in a galaxy cluster consists of dark matter and not baryonic matter. For the Virgo Cluster, the ICM contains roughly 3 × 1014 M☉ while the total mass of the cluster is estimated to be 1.2 × 1015 M☉.
Although the ICM on the whole contains the bulk of a cluster's baryons, it is not very dense, with typical values of 10−3 particles per cubic centimeter. The mean free path of the particles is roughly 1016 m, or about one lightyear. The density of the ICM rises towards the centre of the cluster with a relatively strong peak. In addition, the temperature of the ICM typically drops to 1/2 or 1/3 of the outer value in the central regions. Once the density of the plasma reaches a critical value, enough interactions between the ions ensures cooling via X-ray radiation.
Observing the intracluster medium
As the ICM is at such high temperatures, it emits X-ray radiation, mainly by the bremsstrahlung process and X-ray emission lines from the heavy elements. These X-rays can be observed using an X-ray telescope and through analysis of this data, it is possible to determine the physical conditions, including the temperature, density, and metallicity of the plasma.
Measurements of the temperature and density profiles in galaxy clusters allow for a determination of the mass distribution profile of the ICM through hydrostatic equilibrium modeling. The mass distributions determined from these methods reveal masses that far exceed the luminous mass seen and are thus a strong indication of dark matter in galaxy clusters.
Inverse Compton scattering of low energy photons through interactions with the relativistic electrons in the ICM cause distortions in the spectrum of the cosmic microwave background radiation (CMB), known as the Sunyaev–Zel'dovich effect. These temperature distortions in the CMB can be used by telescopes such as the South Pole Telescope to detect dense clusters of galaxies at high redshifts.
In December 2022, the James Webb Space Telescope is reported to be studying the faint light emitted in the intracluster medium. Which a 2018 study found to be an "accurate luminous tracer of dark matter".
Cooling flows
Plasma in regions of the cluster, with a cooling time shorter than the age of the system, should be cooling due to strong X-ray radiation where emission is proportional to the density squared. Since the density of the ICM is highest towards the center of the cluster, the radiative cooling time drops a significant amount. The central cooled gas can no longer support the weight of the external hot gas and the pressure gradient drives what is known as a cooling flow where the hot gas from the external regions flows slowly towards the center of the cluster. This inflow would result in regions of cold gas and thus regions of new star formation. Recently however, with the launch of new X-ray telescopes such as the Chandra X-ray Observatory, images of galaxy clusters with better spatial resolution have been taken. These new images do not indicate signs of new star formation on the order of what was historically predicted, motivating research into the mechanisms that would prevent the central ICM from cooling.
Heating
There are two popular explanations of the mechanisms that prevent the central ICM from cooling: feedback from active galactic nuclei through injection of relativistic jets of plasma and sloshing of the ICM plasma during mergers with subclusters. The relativistic jets of material from active galactic nuclei can be seen in images taken by telescopes with high angular resolution such as the Chandra X-ray Observatory.
| Physical sciences | Basics_2 | Astronomy |
11983318 | https://en.wikipedia.org/wiki/Branches%20of%20science | Branches of science | The branches of science, also referred to as sciences, scientific fields or scientific disciplines, are commonly divided into three major groups:
Formal sciences: the study of formal systems, such as those under the branches of logic and mathematics, which use an a priori, as opposed to empirical, methodology. They study abstract structures described by formal systems.
Natural sciences: the study of natural phenomena (including cosmological, geological, physical, chemical, and biological factors of the universe). Natural science can be divided into two main branches: physical science and life science (or biology).
Social sciences and the humanities: the study of human behavior in its social and cultural aspects.
Scientific knowledge must be grounded in observable phenomena and must be capable of being verified by other researchers working under the same conditions.
Natural, social, and formal science make up the fundamental sciences, which form the basis of interdisciplinarity - and applied sciences such as engineering and medicine. Specialized scientific disciplines that exist in multiple categories may include parts of other scientific disciplines but often possess their own terminologies and expertises.
Formal sciences
The formal sciences are the branches of science that are concerned with formal systems, such as logic, mathematics, theoretical computer science, information theory, systems theory, decision theory, statistics.
Unlike other branches, the formal sciences are not concerned with the validity of theories based on observations in the real world (empirical knowledge), but rather with the properties of formal systems based on definitions and rules. Hence there is disagreement on whether the formal sciences actually constitute as a science. Methods of the formal sciences are, however, essential to the construction and testing of scientific models dealing with observable reality, and major advances in formal sciences have often enabled major advances in the empirical sciences.
Logic
Logic (from Greek: ) is the systematic study of valid rules of inference, i.e. the relations that lead to the acceptance of one proposition (the conclusion) on the basis of a set of other propositions (premises). More broadly, logic is the analysis and appraisal of arguments.
It has traditionally included the classification of arguments; the systematic exposition of the logical forms; the validity and soundness of deductive reasoning; the strength of inductive reasoning; the study of formal proofs and inference (including paradoxes and fallacies); and the study of syntax and semantics.
Historically, logic has been studied in philosophy (since ancient times) and mathematics (since the mid-19th century). More recently, logic has been studied in cognitive science, which draws on computer science, linguistics, philosophy and psychology, among other disciplines.
Data science
Information science
Information science is an academic field which is primarily concerned with analysis, collection, classification, manipulation, storage, retrieval, movement, dissemination, and protection of information. Practitioners within and outside the field study the application and the usage of knowledge in organizations in addition to the interaction between people, organizations, and any existing information systems with the aim of creating, replacing, improving, or understanding the information systems.
Mathematics
Mathematics, in the broadest sense, is just a synonym of formal science; but traditionally mathematics means more specifically the coalition of four areas: arithmetic, algebra, geometry, and analysis, which are, to some degree, the study of quantity, structure, space, and change respectively.
Statistics
Statistics is the study of the collection, organization, and interpretation of data. It deals with all aspects of this, including the planning of data collection in terms of the design of surveys and experiments.
A statistician is someone who is particularly well versed in the ways of thinking necessary for the successful application of statistical analysis. Such people have often gained this experience through working in any of a wide number of fields. There is also a discipline called mathematical statistics, which is concerned with the theoretical basis of the subject.
The word statistics, when referring to the scientific discipline, is singular, as in "Statistics is an art." This should not be confused with the word statistic, referring to a quantity (such as mean or median) calculated from a set of data, whose plural is statistics ("this statistic seems wrong" or "these statistics are misleading").
Systems theory
Systems theory is the transdisciplinary study of systems in general, to elucidate principles that can be applied to all types of systems in all fields of research. The term does not yet have a well-established, precise meaning, but systems theory can reasonably be considered a specialization of systems thinking and a generalization of systems science. The term originates from Bertalanffy's General System Theory (GST) and is used in later efforts in other fields, such as the action theory of Talcott Parsons and the sociological autopoiesis of Niklas Luhmann.
In this context the word systems is used to refer specifically to self-regulating systems, i.e. that are self-correcting through feedback. Self-regulating systems are found in nature, including the physiological systems of the human body, in local and global ecosystems, and climate.
Decision theory
Decision theory (or the theory of choice not to be confused with choice theory) is the study of an agent's choices. Decision theory can be broken into two branches: normative decision theory, which analyzes the outcomes of decisions or determines the optimal decisions given constraints and assumptions, and descriptive decision theory, which analyzes how agents actually make the decisions they do.
Decision theory is closely related to the field of game theory and is an interdisciplinary topic, studied by economists, statisticians, psychologists, biologists, political and other social scientists, philosophers, and computer.
Empirical applications of this rich theory are usually done with the help of statistical and econometric methods.
Theoretical computer science
Theoretical computer science (TCS) is a subset of general computer science and mathematics that focuses on more mathematical topics of computing, and includes the theory of computation.
It is difficult to circumscribe the theoretical areas precisely. The ACM's (Association for Computing Theory) Special Interest Group on Algorithms and Computation Theory (SIGACT) provides the following description:
Natural sciences
Natural science is a branch of science concerned with the description, prediction, and understanding of natural phenomena, based on empirical evidence from observation and experimentation. Mechanisms such as peer review and repeatability of findings are used to try to ensure the validity of scientific advances.
Natural science can be divided into two main branches: life science and physical science. Life science is alternatively known as biology, and physical science is subdivided into branches: physics, chemistry, astronomy and Earth science. These branches of natural science may be further divided into more specialized branches (also known as fields).
Physical science
Physical science is an encompassing term for the branches of natural science that study non-living systems, in contrast to the life sciences. However, the term "physical" creates an unintended, somewhat arbitrary distinction, since many branches of physical science also study biological phenomena. There is a difference between physical science and physics.
Physics
Physics (from ) is a natural science that involves the study of matter and its motion through spacetime, along with related concepts such as energy and force. More broadly, it is the general analysis of nature, conducted in order to understand how the universe behaves.
Physics is one of the oldest academic disciplines, perhaps the oldest through its inclusion of astronomy. Over the last two millennia, physics was a part of natural philosophy along with chemistry, certain branches of mathematics, and biology, but during the Scientific Revolution in the 16th century, the natural sciences emerged as unique research programs in their own right. Certain research areas are interdisciplinary, such as biophysics and quantum chemistry, which means that the boundaries of physics are not rigidly defined. In the nineteenth and twentieth centuries physicalism emerged as a major unifying feature of the philosophy of science as physics provides fundamental explanations for every observed natural phenomenon. New ideas in physics often explain the fundamental mechanisms of other sciences, while opening to new research areas in mathematics and philosophy.
Chemistry
Chemistry (the etymology of the word has been much disputed) is the science of matter and the changes it undergoes. The science of matter is also addressed by physics, but while physics takes a more general and fundamental approach, chemistry is more specialized, being concerned by the composition, behavior (or reaction), structure, and properties of matter, as well as the changes it undergoes during chemical reactions. It is a physical science which studies various substances, atoms, molecules, and matter (especially carbon based). Example sub-disciplines of chemistry include: biochemistry, the study of substances found in biological organisms; physical chemistry, the study of chemical processes using physical concepts such as thermodynamics and quantum mechanics; and analytical chemistry, the analysis of material samples to gain an understanding of their chemical composition and structure. Many more specialized disciplines have emerged in recent years, e.g. neurochemistry the chemical study of the nervous system.
Earth science
Earth science (also known as geoscience, the geosciences or the Earth sciences) is an all-embracing term for the sciences related to the planet Earth. It is arguably a special case in planetary science, the Earth being the only known life-bearing planet. There are both reductionist and holistic approaches to Earth sciences. The formal discipline of Earth sciences may include the study of the atmosphere, hydrosphere, lithosphere, and biosphere, as well as the solid earth. Typically Earth scientists will use tools from physics, chemistry, biology, geography, chronology and mathematics to build a quantitative understanding of how the Earth system works, and how it evolved to its current state.
Geology
Geology (from the Ancient Greek γῆ, gē ("earth") and -λoγία, -logia, ("study of", "discourse")) is an Earth science concerned with the solid Earth, the rocks of which it is composed, and the processes by which they change over time. Geology can also include the study of the solid features of any terrestrial planet or natural satellite such as Mars or the Moon. Modern geology significantly overlaps all other Earth sciences, including hydrology and the atmospheric sciences, and so is treated as one major aspect of integrated Earth system science and planetary science.
Oceanography
Oceanography, or marine science, is the branch of Earth science that studies the ocean. It covers a wide range of topics, including marine organisms and ecosystem dynamics; ocean currents, waves, and geophysical fluid dynamics; plate tectonics and the geology of the seafloor; and fluxes of various chemical substances and physical properties within the ocean and across its boundaries. These diverse topics reflect multiple disciplines that oceanographers blend to further knowledge of the world ocean and understanding of processes within it: biology, chemistry, geology, meteorology, and physics as well as geography.
Meteorology
Meteorology is the interdisciplinary scientific study of the atmosphere. Studies in the field stretch back millennia, though significant progress in meteorology did not occur until the 17th century. The 19th century saw breakthroughs occur after observing networks developed across several countries. After the development of the computer in the latter half of the 20th century, breakthroughs in weather forecasting were achieved.
Astronomy
Space science is the study of everything in outer space. This has sometimes been called astronomy, but recently astronomy has come to be regarded as a division of broader space science, which has grown to include other related fields, such as studying issues related to space travel and space exploration (including space medicine), space archaeology and science performed in outer space (see space research).
Biological science
Life science, also known as biology, is the natural science that studies life such as microorganisms, plants, and animals including human beings, – including their physical structure, chemical processes, molecular interactions, physiological mechanisms, development, and evolution. Despite the complexity of the science, certain unifying concepts consolidate it into a single, coherent field. Biology recognizes the cell as the basic unit of life, genes as the basic unit of heredity, and evolution as the engine that propels the creation and extinction of species. Living organisms are open systems that survive by transforming energy and decreasing their local entropy to maintain a stable and vital condition defined as homeostasis.
Biochemistry
Biochemistry, or biological chemistry, is the study of chemical processes within and relating to living organisms. It is a sub-discipline of both biology and chemistry, and from a reductionist point of view it is fundamental in biology. Biochemistry is closely related to molecular biology, cell biology, genetics, and physiology.
Microbiology
Microbiology is the study of microorganisms, those being unicellular (single cell), multicellular (cell colony), or acellular (lacking cells). Microbiology encompasses numerous sub-disciplines including virology, bacteriology, protistology, mycology, immunology and parasitology.
Botany
Botany, also called plant science(s), plant biology or phytology, is the science of plant life and a branch of biology. Traditionally, botany has also included the study of fungi and algae by mycologists and phycologists respectively, with the study of these three groups of organisms remaining within the sphere of interest of the International Botanical Congress. Nowadays, botanists (in the strict sense) study approximately 410,000 species of land plants of which some 391,000 species are vascular plants (including approximately 369,000 species of flowering plants), and approximately 20,000 are bryophytes.
Zoology
Zoology () is the branch of biology that studies the animal kingdom, including the structure, embryology, evolution, classification, habits, and distribution of all animals, both living and extinct, and how they interact with their ecosystems. The term is derived from Ancient Greek ζῷον, zōion, i.e. "animal" and λόγος, logos, i.e. "knowledge, study". Some branches of zoology include: anthrozoology, arachnology, archaeozoology, cetology, embryology, entomology, helminthology, herpetology, histology, ichthyology, malacology, mammalogy, morphology, nematology, ornithology, palaeozoology, pathology, primatology, protozoology, taxonomy, and zoogeography.
Ecology
Ecology (from , "house", or "environment"; , "study of") is a branch of biology concerning interactions among organisms and their biophysical environment, which includes both biotic and abiotic components. Topics of interest include the biodiversity, distribution, biomass, and populations of organisms, as well as cooperation and competition within and between species. Ecosystems are dynamically interacting systems of organisms, the communities they make up, and the non-living components of their environment. Ecosystem processes, such as primary production, pedogenesis, nutrient cycling, and niche construction, regulate the flux of energy and matter through an environment. Organisms with specific life history traits sustain these processes.
Social sciences
Social science is the branch of science devoted to the study of societies and the relationships among individuals within those societies. The term was formerly used to refer to the field of sociology, the original "science of society", established in the 19th century. In addition to sociology, it now encompasses a wide array of academic disciplines, including anthropology, archaeology, economics, education, history, human geography, law, linguistics, political science, and psychology.
Positivist social scientists use methods resembling those of the natural sciences as tools for understanding society, and so define science in its stricter modern sense. Interpretivist social scientists, by contrast, may use social critique or symbolic interpretation rather than constructing empirically falsifiable theories. In modern academic practice, researchers are often eclectic, using multiple methodologies (for instance, by combining both quantitative and qualitative research). The term "social research" has also acquired a degree of autonomy as practitioners from various disciplines share in its aims and methods.
Applied sciences
Applied science is the use of existing scientific knowledge to achieve practical goals, like technology or inventions.
Within natural science, disciplines that are basic science develop basic information to explain and perhaps predict phenomena in the natural world. Applied science is the use of scientific processes and knowledge as the means to achieve a particularly practical or useful result. This includes a broad range of applied science-related fields, including agricultural science, engineering and medicine.
Applied science can also apply formal science, such as statistics and probability theory, as in epidemiology. Genetic epidemiology is an applied science applying both biological and statistical methods.
Relationships between the branches
The relationships between the branches of science are summarized by the table
Visualizations and metascience
Metascience refers to or includes a field of science that is about science itself. OpenAlex and Scholia can be used to visualize and explore scientific fields and research topics.
| Physical sciences | Science basics | Basics and measurement |
5617916 | https://en.wikipedia.org/wiki/Propylene%20carbonate | Propylene carbonate | Propylene carbonate (often abbreviated PC) is an organic compound with the formula C4H6O3. It is a cyclic carbonate ester derived from propylene glycol. This colorless and odorless liquid is useful as a polar, aprotic solvent. Propylene carbonate is chiral, but is used as the racemic mixture in most contexts.
Preparation
Although many organic carbonates are produced using phosgene, propylene and ethylene carbonates are exceptions. They are mainly prepared by the carbonation of the epoxides (epoxypropane, or propylene oxide here):
CH3CHCH2O + CO2 → CH3C2H3O2CO
The corresponding reaction of 1,2-propanediol with phosgene is complex, yielding not only propylene carbonate but also oligomeric products.
Propylene carbonate can also be synthesized from urea and propylene glycol over zinc acetate.
Applications
As a solvent
Propylene carbonate is used as a polar, aprotic solvent. It has a high molecular dipole moment (4.9 D), considerably higher than those of acetone (2.91 D) and ethyl acetate (1.78 D). It is possible, for example, to obtain potassium, sodium, and other alkali metals by electrolysis of their chlorides and other salts dissolved in propylene carbonate.
Electrolyte
Due to its high relative permittivity (dielectric constant) of 64, it is frequently used as a high-permittivity component of electrolytes in lithium batteries, usually together with a low-viscosity solvent (e.g. dimethoxyethane). Its high polarity allows it to create an effective solvation shell around lithium ions, thereby creating a conductive electrolyte. However, it is not used in lithium-ion batteries due to its destructive effect on graphite.
Other
Propylene carbonate can also be found in some adhesives, paint strippers, and in cosmetics. It is also used as plasticizer. Propylene carbonate is also used as a solvent for removal of CO2 from natural gas and synthesis gas where H2S is not also present. This use was developed by El Paso Natural Gas Company and Fluor Corporation in the 1950s for use at the Terrell County Gas Plant in West Texas, now owned by Occidental Petroleum.
Propylene carbonate product may be converted to other carbonate esters by transesterification as well (see Carbonate ester#Carbonate transesterification).
In electrospray ionization mass spectrometry, propylene carbonate is doped into low surface tension solutions to increase analyte charging.
In Grignard reaction propylene carbonate (or most other carbonate esters) might be used to create tertiary alcohols.
Safety
Clinical studies indicate that propylene carbonate does not cause skin irritation or sensitization when used in cosmetic preparations, whereas moderate skin irritation is observed when used undiluted. No significant toxic effects were observed in rats fed propylene carbonate, exposed to the vapor, or exposed to the undiluted liquid. In the US, propylene carbonate is not regulated as a volatile organic compound (VOC) because it does not contribute significantly to the formation of smog and because its vapor is not known or suspected to cause cancer or other toxic effects.
| Physical sciences | Esters and ethers | Chemistry |
5622659 | https://en.wikipedia.org/wiki/Photometric%20system | Photometric system | In astronomy, a photometric system is a set of well-defined passbands (or optical filters), with a known sensitivity to incident radiation. The sensitivity usually depends on the optical system, detectors and filters used. For each photometric system a set of primary standard stars is provided.
A commonly adopted standardized photometric system is the Johnson-Morgan or UBV photometric system (1953). At present, there are more than 200 photometric systems.
Photometric systems are usually characterized according to the widths of their passbands:
broadband (passbands wider than 30 nm, of which the most widely used is Johnson-Morgan UBV system)
intermediate band (passbands between 10 and 30 nm wide)
narrow band (passbands less than 10 nm wide)
Photometric letters
Each letter designates a section of light of the electromagnetic spectrum; these cover well the consecutive major groups, near-ultraviolet (NUV), visible light (centered on the V band), near-infrared (NIR) and part of mid-infrared (MIR). The letters are not standards, but are recognized by common agreement among astronomers and astrophysicists.
The use of U,B,V,R,I bands dates from the 1950s, being single-letter abbreviations.
With the advent of infrared detectors in the next decade, the J to N bands were labelled following on from near-infrared's closest-to-red band, I.
Later the H band was inserted, then Z in the 1990s and finally Y, without changing earlier definitions. Hence, H is out of alphabetical order from its neighbours, while Z,Y are reversed from the alphabetical – higher-wavelength – sub-series which dominates current photometric bands.
Note: colors are only approximate and based on wavelength to sRGB representation (when possible).
Combinations of these letters are frequently used; for example the combination JHK has been used more or less as a synonym of "near-infrared", and appears in the title of many papers.
Filters used
The filters currently being used by other telescopes or organizations.
Units of measurements:
Å = Ångström
nm = nanometre
μm = micrometre
Note: colors are only approximate and based on wavelength to sRGB representation (when possible).
| Physical sciences | Basics | Astronomy |
5622894 | https://en.wikipedia.org/wiki/Rh%20blood%20group%20system | Rh blood group system | The Rh blood group system is a human blood group system. It contains proteins on the surface of red blood cells. After the ABO blood group system, it is the most likely to be involved in transfusion reactions. The Rh blood group system consisted of 49 defined blood group antigens in 2005. there are over 50 antigens among which the five antigens D, C, c, E, and e are among the most prominent. There is no d antigen. Rh(D) status of an individual is normally described with a positive (+) or negative (−) suffix after the ABO type (e.g., someone who is A+ has the A antigen and Rh(D) antigen, whereas someone who is A− has the A antigen but lacks the Rh(D) antigen). The terms Rh factor, Rh positive, and Rh negative refer to the Rh(D) antigen only. Antibodies to Rh antigens can be involved in hemolytic transfusion reactions and antibodies to the Rh(D) and Rh antigens confer significant risk of hemolytic disease of the newborn.
Nomenclature
The Rh blood group system has two sets of nomenclature: one developed by Ronald Fisher and R. R. Race, the other by .The two systems reflect different theories of inheritance. The Fisher–Race system, which is currently more common, uses the CDE nomenclature. This system is based on the theory that a separate gene controls the product of each corresponding antigen (e.g., a "D gene" produces D antigen, and so on). However, the d gene was hypothetical, not actual.
The Wiener system uses the Rh–Hr nomenclature. This system is based on the theory that there is one gene at a single locus on each of the two copies of chromosome 1, each contributing to production of multiple antigens. In this theory, a gene R1 is supposed to give rise to the "blood factors" Rh0, rh′, and rh″ (corresponding to modern nomenclature of the D, C, and E antigens) and the gene r to produce hr′ and hr″ (corresponding to modern nomenclature of the c and e antigens).
Notations of the two theories are used interchangeably in blood banking (e.g., Rho(D) meaning RhD positive). Wiener's notation is more complex and cumbersome for routine use. Because it is simpler to explain, the Fisher–Race theory has become more widely used.
DNA testing has shown that both are partially correct: There are in fact two linked genes, the RHD gene which produces a single immune specificity (anti-D) and the RHCE gene with multiple specificities (anti-C, anti-c, anti-E, anti-e). Thus, Wiener's postulate that a gene could have multiple specificities (something many did not give credence to originally) has been proved to be correct. On the other hand, Wiener's theory that there is only one gene has proved to be incorrect, as has the Fisher–Race theory that there are three genes, rather than the two. The CDE notation used in the Fisher–Race nomenclature is sometimes rearranged to DCE to more accurately represent the co-location of the C and E encoding on the RhCE gene, and to make interpretation easier.
Antigens
The proteins which carry the Rh antigens are transmembrane proteins, whose structure suggests that they are ion channels. The main antigens are D, C, E, c and e, which are encoded by two adjacent gene loci, the RHD gene which encodes the RhD protein with the D antigen (and variants) and the RHCE gene which encodes the RhCE protein with the C, E, c and e antigens (and variants). There is no d antigen. Lowercase "d" indicates the absence of the D antigen (the gene is usually deleted or otherwise nonfunctional).
Rh phenotypes are readily identified through the presence or absence of the Rh surface antigens. As can be seen in the table below, most of the Rh phenotypes can be produced by several different Rh genotypes. The exact genotype of any individual can only be identified by DNA analysis. Regarding patient treatment, only the phenotype is usually of any clinical significance to ensure a patient is not exposed to an antigen they are likely to develop antibodies against. A probable genotype may be speculated on, based upon the statistical distributions of genotypes in the patient's place of origin.
R0 (cDe or Dce) is today most common in Africa. The allele was thus often assumed in early blood group analyses to have been typical of populations on the continent, particularly in areas below the Sahara. Ottensooser et al. (1963) suggested that high R0 frequencies were likely characteristic of the ancient Judean Jews, who had emigrated from Egypt prior to their dispersal throughout the Mediterranean Basin and Europe on the basis of high R0 percentages among Sephardi and Ashkenazi Jews compared to native European populations and the relative genetic isolation of Ashkenazim. However, more recent studies have found R0 frequencies as low as 24.3% among some Afroasiatic-speaking groups in the Horn of Africa, as well as higher R0 frequencies among certain other Afroasiatic speakers in North Africa (37.3%) and among some Palestinians in the Levant (30.4%). On the contrary, at a frequency of 47.2% of the population of Basque country having the lack of the D antigen, these people display the highest frequency of the Rh negative phenotype.
• Figures taken from a study performed in 1948 on a sample of 2000 people in the United Kingdom.
Rh antibodies
Rh antibodies are Immunoglobulin G (IgG) antibodies which are acquired through exposure to Rh-positive blood (generally either through pregnancy or transfusion of blood products). The D antigen is the most immunogenic of all the non-ABO antigens. Approximately 80% of individuals who are D-negative and exposed to a single D-positive unit will produce an anti-D antibody. The percentage of alloimmunization is significantly reduced in patients who are actively exsanguinating.
All Rh antibodies except D display dosage (antibody reacts more strongly with red cells homozygous for an antigen than cells heterozygous for the antigen (EE stronger reaction vs Ee)).
If anti-E is detected, the presence of anti-c should be strongly suspected (due to combined genetic inheritance). It is therefore common to select c-negative and E-negative blood for transfusion patients who have an anti-E and lack the c antigen (in general, a patient will not produce antibodies against their own antigens). Anti-c is a common cause of delayed hemolytic transfusion reactions.
Hemolytic disease of the newborn
The hemolytic condition occurs when there is an incompatibility between the blood types of the mother and fetus. There is also potential incompatibility if the mother is Rh negative, and the father is positive. When the mother conceives for the first time, with a positive child, she will become extremely sensitive. When any incompatibility is detected when she conceives the second time in less than two years then, the mother often receives an injection at 28 weeks' gestation and at birth to avoid the development of antibodies towards the fetus. If not given, then the baby will be dead and must be aborted. These terms do not indicate which specific antigen-antibody incompatibility is implicated. The disorder in the fetus due to Rh D incompatibility is known as erythroblastosis fetalis.
Hemolytic comes from two words: "hema" (blood) and "lysis" (solution) or breaking down of red blood cells
Erythroblastosis refers to the making of immature red blood cells
Fetalis refers to the fetus.
When the condition is caused by the Rh D antigen-antibody incompatibility, it is called Rh D Hemolytic disease of the newborn or Rh disease. Here, sensitization to Rh D antigens (usually by feto-maternal transfusion during pregnancy) may lead to the production of maternal IgG anti-D antibodies which can pass through the placenta. This is of particular importance to D negative females at or below childbearing age, because any subsequent pregnancy may be affected by the Rh D hemolytic disease of the newborn if the baby is D positive. The vast majority of Rh disease is preventable in modern antenatal care by injections of IgG anti-D antibodies (Rho(D) Immune Globulin). The incidence of Rh disease is mathematically related to the frequency of D negative individuals in a population, so Rh disease is rare in old-stock populations of Africa and the eastern half of Asia, and the Indigenous peoples of Oceania and the Americas, but more common in other genetic groups, most especially Western Europeans, but also other West Eurasians, and to a lesser degree, native Siberians, as well as those of mixed-race with a significant or dominant descent from those (e.g. the vast majority of Latin Americans and Central Asians).
Symptoms and signs in the fetus:
Enlarged liver, spleen, or heart and fluid buildup in the fetus' abdomen seen via ultrasound.
Symptoms and signs in the newborn:
Anemia that creates the newborn's pallor (pale appearance).
Jaundice or yellow discoloration of the newborn's skin, sclera or mucous membrane. This may be evident right after birth or after 24–48 hours after birth. This is caused by bilirubin (one of the end products of red blood cell destruction).
Enlargement of the newborn's liver and spleen.
The newborn may have severe edema of the entire body.
Dyspnea (difficulty breathing)
Other animals with Rh-like antigens causing hemolytic disease of the newborn
Rh disease only occurs in human fetuses, however a similar disease called Neonatal isoerythrolysis (NI) can be observed in animal species of newborn horses, mules, pigs, cats, cattle, and dogs. What differs between Rh disease and NI is the pathogenesis of hemolysis between human fetuses and the animal species. With human mothers, the maternal antibodies are formed from the sensitization of foreign antigens of her unborn fetus’s red blood cells passing through the placenta causing hemolysis before birth. With other animals, however, these maternal antibodies are not passed through the placenta, but through colostrum. The newborn animal is without NI but soon develops hemolytic anemia after initial ingestion of its mother’s colostrum that contain antibodies that can be absorbed through the newborn’s intestines and are incompatible to its red blood cell antigen. After 48 hours of birth, the newborn may be allowed to nurse from its mother as her antibodies can no longer be absorbed through the neonate’s intestines. Because the most active newborn animals consume the most colostrum, they may be the ones who are most affected by the blood incompatibility of antigen and antibody.
Population data
According to a comprehensive study, the worldwide frequency of Rh-positive and Rh-negative blood types is approximately 94% and 6%, respectively. The same study concluded that the share of the population with Rh-negative blood type is set to fall further in the future primarily due to low population growth in Europe. The frequency of Rh factor blood types and the RhD neg allele gene differs in various populations.
Genetics
The D antigen is inherited as one gene (RHD) (on the short arm of the first chromosome, p36.13–p34.3) with various alleles. Typically, Rhesus positive people have an intact RHD gene while negative people lack the gene (or have mutations in it). However, there are exceptions: for instance, Japanese and black Africans may have an intact gene that is not expressed or only at very low levels. The gene codes for the RhD protein on the red blood cell membrane. D− individuals who lack a functional RHD gene do not produce the D antigen and may be immunized by D+ blood.
The D antigen is a dominant trait. If both of a child's parents are Rh negative, the child will definitely be Rh negative. Otherwise, the child may be Rh positive or Rh negative, depending on the parents' specific genotypes.
The epitopes for the next 4 most common Rh antigens, C, c, E and e are expressed on the highly similar RhCE protein that is genetically encoded in the RHCE gene, also found on chromosome 1. It has been shown that the RHD gene arose by duplication of the RHCE gene during primate evolution. Mice have just one RH gene.
The RHAG gene, which is responsible for encoding Rh-associated glycoprotein (RhAG), is found on chromosome 6a.
The polypeptides produced from the RHD and RHCE genes form a complex on the red blood cell membrane with the Rh-associated glycoprotein.
Function
On the basis of structural homology it has been proposed that the product of RHD gene, the RhD protein, is a membrane transport protein of uncertain specificity (CO2 or NH3) and unknown physiological role. The three-dimensional structure of the related RHCG protein and biochemical analysis of the RhD protein complex indicates that the RhD protein is one of three subunits of an ammonia transporter. Three recent studies have reported a protective effect of the RhD-positive phenotype, especially RhD heterozygosity, against the negative effect of latent toxoplasmosis on psychomotor performance in infected subjects. RhD-negative compared to RhD-positive subjects without anamnestic titres of anti-Toxoplasma antibodies have shorter reaction times in tests of simple reaction times. And conversely, RhD-negative subjects with anamnestic titres (i.e. with latent toxoplasmosis) exhibited much longer reaction times than their RhD-positive counterparts. The published data suggested that only the protection of RhD-positive heterozygotes was long term in nature; the protection of RhD-positive homozygotes decreased with duration of the infection while the performance of RhD-negative homozygotes decreased immediately after the infection. The overall change in reaction times was always larger in the RhD-negative group than in the RhD-positive.
Non-human Rh proteins
Rh-like proteins can be found even in species other than vertebrates (which have red blood cells) – worms, bacteria, and algae. All these Rh proteins have the same biochemical function of transporting , differing slightly in their amino acid sequences. The Rh family as a whole is related to ammonia transporters (Amt). In C. elegans worms, disruption of the Rh1 gene causes growth defects under high levels. Chlamydomonas reinhardtii algae fail to grow rapidly if its Rh gene is knocked down. Although in vitro evidence shows the Rh complex is capable of moving ammonia, its disruption does not cause growth defects under modified ammonia levels.
RHD polymorphism
Origin of RHD polymorphism
For a long time, the origin of RHD polymorphism was an evolutionary enigma. Before the advent of modern medicine, the carriers of the rarer allele (e.g. RhD-negative women in a population of RhD positives or RhD-positive men in a population of RhD negatives) were at a disadvantage as some of their children (RhD-positive children born to preimmunised RhD-negative mothers) were at a higher risk of fetal or newborn death or health impairment from hemolytic disease.
Natural selection aside, the RHD-RHCE region is structurally predisposed to many mutations seen in humans, since the pair arose by gene duplication and remain similar enough for unequal crossing over to occur. In addition to the case where D is deleted, crossover can also produce a single gene mixing exons from both RHD and RHCE, forming the majority of partial D types.
Weak D
In serologic testing, D positive blood is easily identified. Units that are D negative are often retested to rule out a weaker reaction. This was previously referred to as Du, which has been replaced. By definition, weak D phenotype is characterized by negative reaction with anti-D reagent at immediate spin (IS), negative reaction after 37 °C incubation, and positive reaction at anti-human globulin (AHG) phase. Weak D phenotype can occur in several ways. In some cases, this phenotype occurs because of an altered surface protein that is more common in people of European descent. An inheritable form also occurs, as a result of a weakened form of the R0 gene. Weak D may also occur as "C in trans", whereby a C gene is present on the opposite chromosome to a D gene (as in the combination R0r', or "Dce/dCe"). The testing is difficult, since using different anti-D reagents, especially the older polyclonal reagents, may give different results.
The practical implication of this is that people with this sub-phenotype will have a product labeled as "D positive" when donating blood. When receiving blood, they are sometimes typed as a "D negative", though this is the subject of some debate. Most "Weak D" patients can receive "D positive" blood without complications. However, it is important to correctly identify the ones that have to be considered D+ or D−. This is important, since most blood banks have a limited supply of "D negative" blood and the correct transfusion is clinically relevant. In this respect, genotyping of blood groups has much simplified this detection of the various variants in the Rh blood group system.
Partial D
It is important to differentiate weak D (due to a quantitative difference in the D antigen) from partial D (due to a qualitative difference in the D antigen). Simply put, the weak D phenotype is due to a reduced number of D antigens on a red blood cell. In contrast, the partial D phenotype is due to an alteration in D-epitopes. Thus, in partial D, the number of D antigens is not reduced but the protein structure is altered. These individuals, if alloimmunized to D, can produce an anti-D antibody. Therefore, partial D patients who are donating blood should be labeled as D-positive but, if receiving blood, they should be labeled as D-negative and receive D-negative units.
In the past, partial D was called 'D mosaic' or 'D variant.' Different partial D phenotypes are defined by different D epitopes on the outer surface of the red blood cell membrane. More than 30 different partial D phenotypes have been described.
Rhnull phenotype
Rhnull individuals have no Rh antigens (no Rh or RhAG) on their red blood cells. This rare condition has been called "Golden Blood". As a consequence of Rh antigen absence, Rhnull red blood cells also lack LW and Fy5 and show weak expression of S, s, and U antigens.
Red blood cells lacking Rh/RhAG proteins have structural abnormalities (such as stomatocytosis) and cell membrane defects that can result in hemolytic anemia.
The first Rhnull blood was discovered in an Aboriginal Australian woman, in 1961. Only 43 individuals have been reported to have it worldwide. Only nine active donors have been reported. Its properties make it attractive in numerous medical applications, but scarcity makes it expensive to transport and acquire.
Other Rh group antigens
As of 2023, over 50 antigens have been described in the Rh group system; among those described here, the D, C, c, E and e antigens are the most important. The others are much less frequently encountered or are rarely clinically significant. Each is given a number, though the highest assigned number (CEVF or RH61 according to the ISBT terminology) is not an accurate reflection of the antigens encountered since many (e.g. Rh38) have been combined, reassigned to other groups, or otherwise removed.
Some of the other Rh "antigens" are f ("ce", RH6), Ce (RH7), Cw (RH8), Cx (RH9), V (RH10), Ew (RH11), G (RH12), Tar (RH40), VS (RH20), Dw (RH23), and CE (RH22). Some of these groups, including f, Ce and CE, describe grouping of some existing groups. Others, like V, describe an epitope created by some other mutation on the RHD and RHCE genes. V in particular is caused by a mutation on RHCE.
History
The term "Rh" was originally an abbreviation of "Rhesus factor". It was discovered in 1939 by Karl Landsteiner and Alexander S. Wiener, who, at the time, believed it to be a similar antigen found in rhesus macaque red blood cells. It was subsequently discovered that the human factor is not identical to the rhesus monkey factor, but by then, "Rhesus Group" and like terms were already in widespread, worldwide use. Thus, notwithstanding it is a misnomer, the term survives (e.g., rhesus blood group system and the obsolete terms rhesus factor, rhesus positive, and rhesus negative – all three of which actually refer specifically and only to the Rh D factor and are thus misleading when unmodified). Contemporary practice is to use "Rh" as a term of art instead of "Rhesus" (e.g., "Rh Group", "Rh factors", "Rh D", etc.).
The significance of their discovery was not immediately apparent and was only realized in 1940, after subsequent findings by Philip Levine and Rufus Stetson. The serum that led to the discovery was produced by immunizing rabbits with red blood cells from a rhesus macaque. The antigen that induced this immunization was designated by them as Rh factor to indicate that rhesus blood had been used for the production of the serum.
In 1939, Phillip Levine and Rufus Stetson published in a first case report the clinical consequences of non-recognized Rh factor, hemolytic transfusion reaction, and hemolytic disease of the newborn in its most severe form. It was recognized that the serum of the reported woman agglutinated with red blood cells of about 80% of the people although the then known blood groups, in particular ABO were matched. No name was given to this agglutinin when described. In 1940, Landsteiner and Wiener made the connection to their earlier discovery, reporting a serum that also reacted with about 85% of different human red blood cells.
In 1941, Group O: a patient in Irvington, New Jersey, US, delivered a normal infant in 1931; this pregnancy was followed by a long period of sterility. The second pregnancy (April, 1941) resulted in an infant with icterus gravis. In May 1941, the third anti-Rh serum (M.S.) of Group O became available.
Based on the serologic similarities, 'Rh factor' was later also used for antigens, and anti-Rh for antibodies, found in humans such as those previously described by Levine and Stetson. Although differences between these two sera were shown already in 1942 and clearly demonstrated in 1963, the already widely used term "Rh" was kept for the clinically described human antibodies which are different from the ones related to the rhesus monkey. This real factor found in rhesus macaque was classified in the Landsteiner-Weiner antigen system (antigen LW, antibody anti-LW) in honor of the discoverers.
It was recognized that the Rh factor was just one in a system of various antigens. Based on different models of genetic inheritance, two different terminologies were developed; both of them are still in use.
The clinical significance of this highly immunizing D antigen (i.e., Rh factor) was soon realized. Some keystones were to recognize its importance for blood transfusion (including reliable diagnostic tests), hemolytic disease of the newborn (including exchange transfusion), and very importantly the prevention of it by screening and prophylaxis.
The discovery of cell-free fetal DNA in maternal circulation by Holzgrieve et al. led to the noninvasive genotyping of fetal Rh genes in many countries.
| Biology and health sciences | Human anatomy | Health |
11005224 | https://en.wikipedia.org/wiki/Penetrating%20trauma | Penetrating trauma | Penetrating trauma is an open wound injury that occurs when an object pierces the skin and enters a tissue of the body, creating a deep but relatively narrow entry wound. In contrast, a blunt or non-penetrating trauma may have some deep damage, but the overlying skin is not necessarily broken and the wound is still closed to the outside environment. The penetrating object may remain in the tissues, come back out the path it entered, or pass through the full thickness of the tissues and exit from another area.
A penetrating injury in which an object enters the body or a structure and passes all the way through an exit wound is called a perforating trauma, while the term penetrating trauma implies that the object does not perforate wholly through. In gunshot wounds, perforating trauma is associated with an entrance wound and an often larger exit wound.
Penetrating trauma can be caused by a foreign object or by fragments of a broken bone. Usually occurring in violent crime or armed combat, penetrating injuries are commonly caused by gunshots and stabbings.
Penetrating trauma can be serious because it can damage internal organs and presents a risk of shock and infection. The severity of the injury varies widely depending on the body parts involved, the characteristics of the penetrating object, and the amount of energy transmitted to the tissues. Assessment may involve X-rays or CT scans, and treatment may involve surgery, for example to repair damaged structures or to remove foreign objects. Following penetrating trauma, spinal motion restriction is associated with worse outcomes and therefore it should not be done routinely.
Mechanism
As a missile passes through tissue, it decelerates, dissipating and transferring kinetic energy to the tissues. The velocity of the projectile is a more important factor than its mass in determining how much damage is done; kinetic energy increases with the square of the velocity. In addition to injury caused directly by the object that enters the body, penetrating injuries may be associated with secondary injuries, due for example to a blast injury.
The path of a projectile can be estimated by imagining a line from the entrance wound to the exit wound, but the actual trajectory may vary due to ricochet or differences in tissue density. In a cut, the discolouration and the swelling of the skin from a blow happens because of the ruptured blood vessels and escape of blood and fluid and other injuries that interrupt the circulation.
Cavitation
Permanent
Low-velocity items, such as knives and swords, are usually propelled by a person's hand, and usually do damage only to the area that is directly contacted by the object. The space left by tissue that is destroyed by the penetrating object as it passes through forms a cavity; this is called permanent cavitation.
Temporary
High-velocity objects are usually projectiles such as bullets from high-powered rifles, such as assault rifles or sniper rifles. Bullets classed as medium-velocity projectiles include those from handguns, shotguns, and submachine guns. In addition to causing damage to the tissues they contact, medium- and high-velocity projectiles cause a secondary cavitation injury: as the object enters the body, it creates a pressure wave which forces tissue out of the way, creating a cavity which can be much larger than the object itself; this is called "temporary cavitation". The temporary cavity is the radial stretching of tissue around the bullet's wound track, which momentarily leaves an empty space caused by high pressures surrounding the projectile that accelerate material away from its path.
The characteristics of the tissue injured also help determine the severity of the injury; for example, the denser the tissue, the greater the amount of energy transmitted to it. Skin, muscles, and intestines absorb energy and so are resistant to the development of temporary cavitation, while organs such as the liver, spleen, kidney, and brain, which have relatively low tensile strength, are likely to split or shatter because of temporary cavitation. Flexible elastic soft tissues, such as muscle, intestine, skin, and blood vessels, are good energy absorbers and are resistant to tissue stretch. If enough energy is transferred, the liver may disintegrate. Temporary cavitation can be especially damaging when it affects delicate tissues such as the brain, as occurs in penetrating head trauma.
Location
Head
While penetrating head trauma accounts for only a small percentage of all traumatic brain injuries (TBI), it is associated with a high mortality rate, and only a third of people with penetrating head trauma survive long enough to arrive at a hospital. Injuries from firearms are the leading cause of TBI-related deaths. Penetrating head trauma can cause cerebral contusions and lacerations, intracranial hematomas, pseudoaneurysms, and arteriovenous fistulas. The prognosis for penetrating head injuries varies widely.
Penetrating facial trauma can pose a risk to the airway and breathing; airway obstruction can occur later due to swelling or bleeding. Penetrating eye trauma can cause the globe of the eye to rupture or vitreous humor to leak from it, and presents a serious threat to eyesight.
Chest
Most penetrating injuries are chest wounds and have a mortality rate (death rate) of under 10%. Penetrating chest trauma can injure vital organs such as the heart and lungs and can interfere with breathing and circulation. Lung injuries that can be caused by penetrating trauma include pulmonary laceration (a cut or tear) pulmonary contusion (a bruise), hemothorax (an accumulation of blood in the chest cavity outside of the lung), pneumothorax (an accumulation of air in the chest cavity) and hemopneumothorax (accumulation of both blood and air). Sucking chest wounds and tension pneumothorax may result.
Penetrating trauma can also cause injuries to the heart and circulatory system. When the heart is punctured, it may bleed profusely into the chest cavity if the membrane around it (the pericardium) is significantly torn, or it may cause pericardial tamponade if the pericardium is not disrupted. In pericardial tamponade, blood escapes from the heart but is trapped within the pericardium, so pressure builds up between the pericardium and the heart, compressing the latter and interfering with its pumping. Fractures of the ribs commonly produce penetrating chest trauma when sharp bone ends pierce tissues.
Abdomen
Penetrating abdominal trauma (PAT) typically arises from stabbings, ballistic injuries (shootings), or industrial accidents. PAT can be life-threatening because abdominal organs, especially those in the retroperitoneal space, can bleed profusely, and the space can hold a large volume of blood. If the pancreas is injured, it may be further injured by its own secretions, in a process called autodigestion. Injuries of the liver, common because of the size and location of the organ, present a serious risk for shock because the liver tissue is delicate and has a large blood supply and capacity. The intestines, taking a large part of the lower abdomen, are also at risk of perforation.
People with penetrating abdominal trauma may have signs of hypovolemic shock (insufficient blood in the circulatory system) and peritonitis (an inflammation of the peritoneum, the membrane that lines the abdominal cavity). Penetration may abolish or diminish bowel sounds due to bleeding, infection, and irritation, and injuries to arteries may cause bruits (a distinctive sound similar to heart murmurs) to be audible. Percussion of the abdomen may reveal hyperresonance (indicating air in the abdominal cavity) or dullness (indicating a buildup of blood). The abdomen may be distended or tender, signs which indicate an urgent need for surgery.
The standard management of penetrating abdominal trauma was for many years mandatory laparotomy. A greater understanding of mechanisms of injury, outcomes from surgery, improved imaging and interventional radiology has led to more conservative operative strategies being adopted.
Assessment and treatment
Assessment can be difficult because much of the damage is often internal and not visible. The patient is thoroughly examined. X-ray and CT scanning may be used to identify the type and location of potentially lethal injuries. Sometimes before an X-ray is performed on a person with penetrating trauma from a projectile, a paper clip is taped over entry and exit wounds to show their location on the film. The patient is given intravenous fluids to replace lost blood. Surgery may be required; impaled objects are secured into place so that they do not move and cause further injury, and they are removed in an operating room. If the location of the injury is not obvious, a surgical operation called an exploratory laparotomy may be required to look for internal damage to the organs in the abdomen. Foreign bodies such as bullets may be removed, but they may also be left in place if the surgery necessary to get them out would cause more damage than would leaving them. Wounds are debrided to remove tissue that cannot survive and other material that presents risk for infection.
Negative pressure wound therapy is no more effective in preventing wound infection than standard care when used on open traumatic wounds.
History
Before the 17th century, medical practitioners poured hot oil into wounds in order to cauterize damaged blood vessels, but the French surgeon Ambroise Paré challenged the use of this method in 1545. Paré was the first to propose controlling bleeding using ligature.
During the American Civil War, chloroform was used during surgery to reduce pain and allow more time for operations. Due in part to the lack of sterile technique in hospitals, infection was the leading cause of death for wounded soldiers.
In World War I, doctors began replacing patients' lost fluid with salt solutions. With World War II came the idea of blood banking, having quantities of donated blood available to replace lost fluids. The use of antibiotics also came into practice in World War II.
| Biology and health sciences | Types | Health |
11007302 | https://en.wikipedia.org/wiki/Halo%20orbit | Halo orbit | A halo orbit is a periodic, three-dimensional orbit associated with one of the L1, L2 or L3 Lagrange points in the three-body problem of orbital mechanics. Although a Lagrange point is just a point in empty space, its peculiar characteristic is that it can be orbited by a Lissajous orbit or by a halo orbit. These can be thought of as resulting from an interaction between the gravitational pull of the two planetary bodies and the Coriolis and centrifugal force on a spacecraft. Halo orbits exist in any three-body system, e.g., a Sun–Earth–orbiting satellite system or an Earth–Moon–orbiting satellite system. Continuous "families" of both northern and southern halo orbits exist at each Lagrange point. Because halo orbits tend to be unstable, station-keeping using thrusters may be required to keep a satellite on the orbit.
Most satellites in halo orbit serve scientific purposes, for example space telescopes.
Definition and history
Robert W. Farquhar first used the name "halo" in 1966 for orbits around L which were made periodic using thrusters. Farquhar advocated using spacecraft in such an orbit beyond the Moon (Earth–Moon ) as a communications relay station for an Apollo mission to the far side of the Moon. A spacecraft in such an orbit would be in continuous view of both the Earth and the far side of the Moon, whereas a Lissajous orbit would sometimes make the spacecraft go behind the Moon. In the end, no relay satellite was launched for Apollo, since all landings were on the near side of the Moon.
In 1973 Farquhar and Ahmed Kamel found that when the in-plane amplitude of a Lissajous orbit was large enough there would be a corresponding out-of-plane amplitude that would have the same period, so the orbit ceased to be a Lissajous orbit and became approximately an ellipse. They used analytical expressions to represent these halo orbits; in 1984, Kathleen Howell showed that more precise trajectories could be computed numerically. Additionally, she found that for most values of the ratio between the masses of the two bodies (such as the Earth and the Moon) there was a range of stable orbits.
The first mission to use a halo orbit was ISEE-3, a joint ESA and NASA spacecraft launched in 1978. It traveled to the Sun–Earth point and remained there for several years. The next mission to use a halo orbit was Solar and Heliospheric Observatory (SOHO), also a joint ESA/NASA mission to study the Sun, which arrived at Sun–Earth in 1996. It used an orbit similar to ISEE-3. Although several other missions since then have traveled to Lagrange points, they (eg. Gaia astrometric space observatory) typically have used the related non-periodic variations called Lissajous orbits rather than an actual halo orbit.
Although halo orbits were well known in the RTBP (Restricted Three Body Problem), it was difficult to obtain Halo orbits for the real Earth-Moon system. Translunar halo orbits were first computed in 1998 by M.A. Andreu, who introduced a new model for the motion of a spacecraft in the Earth-Moon-Sun system, which was called Quasi-Bicircular Problem (QBCP).
In May 2018, Farquhar's original idea was finally realized when China placed the first communications relay satellite, Queqiao, into a halo orbit around the Earth-Moon point. On 3 January 2019, the Chang'e 4 spacecraft landed in the Von Kármán crater on the far side of the Moon, using the Queqiao relay satellite to communicate with the Earth.
The James Webb Space Telescope entered a halo orbit around the Sun-Earth point on 24 January 2022. Euclid entered a similar orbit around this point in August 2023.
India's space agency ISRO launched Aditya-L1 to study the sun from a halo orbit around L point. On 6 January 2024, Aditya-L1 spacecraft, India's first solar mission, has successfully entered its final orbit with a period of approximately 180 days around the first Sun-Earth Lagrangian point (L1), approximately 1.5 million kilometers from Earth.
| Physical sciences | Orbital mechanics | Astronomy |
11008314 | https://en.wikipedia.org/wiki/Evolutionary%20history%20of%20plants | Evolutionary history of plants | The evolution of plants has resulted in a wide range of complexity, from the earliest algal mats of unicellular archaeplastids evolved through endosymbiosis, through multicellular marine and freshwater green algae, to spore-bearing terrestrial bryophytes, lycopods and ferns, and eventually to the complex seed-bearing gymnosperms and angiosperms (flowering plants) of today. While many of the earliest groups continue to thrive, as exemplified by red and green algae in marine environments, more recently derived groups have displaced previously ecologically dominant ones; for example, the ascendance of flowering plants over gymnosperms in terrestrial environments.
There is evidence that cyanobacteria and multicellular thalloid eukaryotes lived in freshwater communities on land as early as 1 billion years ago, and that communities of complex, multicellular photosynthesizing organisms existed on land in the late Precambrian, around .
Evidence of the emergence of embryophyte land plants first occurs in the middle Ordovician (~), and by the middle of the Devonian (~), many of the features recognised in land plants today were present, including roots and leaves. By the late Devonian (~) some free-sporing plants such as Archaeopteris had secondary vascular tissue that produced wood and had formed forests of tall trees. Also by the late Devonian, Elkinsia, an early seed fern, had evolved seeds.
Evolutionary innovation continued throughout the rest of the Phanerozoic eon and still continues today. Most plant groups were relatively unscathed by the Permo-Triassic extinction event, although the structures of communities changed. This may have set the scene for the appearance of the flowering plants in the Triassic (~), and their later diversification in the Cretaceous and Paleogene. The latest major group of plants to evolve were the grasses, which became important in the mid-Paleogene, from around . The grasses, as well as many other groups, evolved new mechanisms of metabolism to survive the low and warm, dry conditions of the tropics over the last .
Colonization of land
Divergence
Land plants evolved from a group of freshwater green algae, perhaps as early as 850 mya, but algae-like plants might have evolved as early as 1 billion years ago. The closest living relatives of land plants are the charophytes, specifically Charales; if modern Charales are similar to the distant ancestors they share with land plants, this means that the land plants evolved from a branched, filamentous alga dwelling in shallow fresh water, perhaps at the edge of seasonally desiccating pools. However, some recent evidence suggests that land plants might have originated from unicellular terrestrial charophytes similar to extant Klebsormidiophyceae. The alga would have had a haplontic life cycle. It would only very briefly have had paired chromosomes (the diploid condition) when the egg and sperm first fused to form a zygote that would have immediately divided by meiosis to produce cells with half the number of unpaired chromosomes (the haploid condition). Co-operative interactions with fungi may have helped early plants adapt to the stresses of the terrestrial realm.
Challenges to land colonization
Plants were not the first photosynthesisers on land. Weathering rates suggest that organisms capable of photosynthesis were already living on the land , and microbial fossils have been found in freshwater lake deposits from , but the carbon isotope record suggests that they were too scarce to impact the atmospheric composition until around . These organisms, although phylogenetically diverse, were probably small and simple, forming little more than an algal scum.
Since lichens initiate the first step in primary ecological succession in contemporary contexts, one hypothesis has been that lichens came on land first and facilitated colonization by plants; however, both molecular phylogenies and the fossil record seem to contradict this.
There are multiple potential reasons for why it took so long for land plants to emerge. It could be that atmospheric 'poisoning' prevented eukaryotes from colonising the land prior to the emergence of land plants, or it could simply have taken a great time for the necessary complexity to evolve. A major challenge to land adaptation would have been the absence of appropriate soil. Throughout the fossil record, soil is preserved, giving information on what early soils were like. Before land plants, the soil on land was poor in resources essential for life like nitrogen and phosphorus and had little capacity for holding water.
Adaptations to land colonization
Evidence of the earliest land plants occurs at about , in lower middle Ordovician rocks from Saudi Arabia and Gondwana in the form of spores known as cryptospores. These spores have walls made of sporopollenin, an extremely decay-resistant material that means they are well-preserved by the fossil record.
These spores were produced either singly (monads), in pairs (dyads) or groups of four (tetrads), and their microstructure resembles that of modern liverwort spores, suggesting they share an equivalent grade of organisation. Their walls contain sporopollenin – further evidence of an embryophytic affinity.
Trilete spores similar to those of vascular plants appear soon afterwards, in Upper Ordovician rocks about 455 million years ago. Depending exactly when the tetrad splits, each of the four spores may bear a "trilete mark", a Y-shape, reflecting the points at which each cell squashed up against its neighbours. However, this requires that the spore walls be sturdy and resistant at an early stage. This resistance is closely associated with having a desiccation-resistant outer wall—a trait only of use when spores must survive out of water. Indeed, even those embryophytes that have returned to the water lack a resistant wall, thus don't bear trilete marks. A close examination of algal spores shows that none have trilete spores, either because their walls are not resistant enough, or, in those rare cases where they are, because the spores disperse before they are compressed enough to develop the mark or do not fit into a tetrahedral tetrad.
The earliest megafossils of land plants were thalloid organisms, which dwelt in fluvial wetlands and are found to have covered most of an early Silurian flood plain. They could only survive when the land was waterlogged. There were also microbial mats.
Once plants had reached the land, there were two approaches to dealing with desiccation. Modern bryophytes either avoid it or give in to it, restricting their ranges to moist settings or drying out and putting their metabolism "on hold" until more water arrives, as in the liverwort genus Targionia. Tracheophytes resist desiccation by controlling the rate of water loss. They all bear a waterproof outer cuticle layer wherever they are exposed to air (as do some bryophytes), to reduce water loss, but since a total covering would cut them off from in the atmosphere tracheophytes use variable openings, the stomata, to regulate the rate of gas exchange. Tracheophytes also developed vascular tissue to aid in the movement of water within the organisms (see below), and moved away from a gametophyte dominated life cycle (see below). Vascular tissue ultimately also facilitated upright growth without the support of water and paved the way for the evolution of larger plants on land.
Consequences
A global glaciation event called Snowball Earth, from around 720-635 mya in the Cryogenian period, is believed to have been at least partially caused by early photosynthetic organisms, which reduced the concentration of carbon dioxide and decreased the greenhouse effect in the atmosphere, leading to an icehouse climate. Based on molecular clock studies of the previous decade or so, a 2022 study observed that the estimated time for the origin of the multicellular streptophytes (all except the unicellular basal clade Mesostigmatophyceae) fell in the cool Cryogenian while that of the subsequent separation of streptophytes fell in the warm Ediacaran, which they interpreted as an indication of selective pressure by the glacial period to the photosynthesizing organisms, a group of which succeeded in surviving in relatively warmer environments that remained habitable, subsequently flourishing in the later Ediacaran and Phanerozoic on land as embryophytes. The study also theorized that the unicellular morphology and other unique features of the Zygnematophyceae may reflect further adaptations to a cold loving life style. The establishment of a land-based flora increased the rate of accumulation of oxygen in the atmosphere, as the land plants produced oxygen as a waste product. When this concentration rose above 13%, around 0.45 billion years ago, wildfires became possible, evident from charcoal in the fossil record. Apart from a controversial gap in the Late Devonian, charcoal has been present ever since.
Charcoalification is an important taphonomic mode. Wildfire or burial in hot volcanic ash drives off the volatile compounds, leaving only a residue of pure carbon. This is not a viable food source for fungi, herbivores or detritovores, so it is prone to preservation. It is also robust and can withstand pressure, displaying exquisite, sometimes sub-cellular, detail in remains.
In addition to the advent of charcoal in the rock record, the terrestrialization of plants has made significant contributions to changes in geology and landscapes. The Ordovician and Silurian show a 1.4 times greater proportion of mudrock in the geologic record than the previous 90% of earth's history and this increase in mudrock is considered to be a result of land plants retaining muds in a terrestrial setting.
Evolution of life cycles
All multicellular plants have a life cycle comprising two generations or phases. The gametophyte phase has a single set of chromosomes (denoted 1n) and produces gametes (sperm and eggs). The sporophyte phase has paired chromosomes (denoted 2n) and produces spores. The gametophyte and sporophyte phases may be homomorphic, appearing identical in some algae, such as Ulva lactuca, but are very different in all modern land plants, a condition known as heteromorphy.
The pattern in plant evolution has been a shift from homomorphy to heteromorphy. The algal ancestors of land plants were almost certainly haplobiontic, being haploid for all their life cycles, with a unicellular zygote providing the 2N stage. All land plants (i.e. embryophytes) are diplobiontic – that is, both the haploid and diploid stages are multicellular. Two trends are apparent: bryophytes (liverworts, mosses and hornworts) have developed the gametophyte as the dominant phase of the life cycle, with the sporophyte becoming almost entirely dependent on it; vascular plants have developed the sporophyte as the dominant phase, with the gametophytes being particularly reduced in the seed plants.
It has been proposed as the basis for the emergence of the diploid phase of the life cycle as the dominant phase that diploidy allows masking of the expression of deleterious mutations through genetic complementation. Thus if one of the parental genomes in the diploid cells contains mutations leading to defects in one or more gene products, these deficiencies could be compensated for by the other parental genome (which nevertheless may have its own defects in other genes). As the diploid phase was becoming predominant, the masking effect likely allowed genome size, and hence information content, to increase without the constraint of having to improve accuracy of replication. The opportunity to increase information content at low cost is advantageous because it permits new adaptations to be encoded. This view has been challenged, with evidence showing that selection is no more effective in the haploid than in the diploid phases of the lifecycle of mosses and angiosperms.
There are two competing theories to explain the appearance of a diplobiontic lifecycle.
The interpolation theory (also known as the antithetic or intercalary theory) holds that the interpolation of a multicellular sporophyte phase between two successive gametophyte generations was an innovation caused by preceding meiosis in a freshly germinated zygote with one or more rounds of mitotic division, thereby producing some diploid multicellular tissue before finally meiosis produced spores. This theory implies that the first sporophytes bore a very different and simpler morphology to the gametophyte they depended on. This seems to fit well with what is known of the bryophytes, in which a vegetative thalloid gametophyte nurtures a simple sporophyte, which consists of little more than an unbranched sporangium on a stalk. Increasing complexity of the ancestrally simple sporophyte, including the eventual acquisition of photosynthetic cells, would free it from its dependence on a gametophyte, as seen in some hornworts (Anthoceros), and eventually result in the sporophyte developing organs and vascular tissue, and becoming the dominant phase, as in the tracheophytes (vascular plants). This theory may be supported by observations that smaller Cooksonia individuals must have been supported by a gametophyte generation. The observed appearance of larger axial sizes, with room for photosynthetic tissue and thus self-sustainability, provides a possible route for the development of a self-sufficient sporophyte phase.
The alternative hypothesis, called the transformation theory (or homologous theory), posits that the sporophyte might have appeared suddenly by delaying the occurrence of meiosis until a fully developed multicellular sporophyte had formed. Since the same genetic material would be employed by both the haploid and diploid phases, they would look the same. This explains the behaviour of some algae, such as Ulva lactuca, which produce alternating phases of identical sporophytes and gametophytes. Subsequent adaption to the desiccating land environment, which makes sexual reproduction difficult, might have resulted in the simplification of the sexually active gametophyte, and elaboration of the sporophyte phase to better disperse the waterproof spores. The tissue of sporophytes and gametophytes of vascular plants such as Rhynia preserved in the Rhynie chert is of similar complexity, which is taken to support this hypothesis. By contrast, modern vascular plants, with the exception of Psilotum, have heteromorphic sporophytes and gametophytes in which the gametophytes rarely have any vascular tissue.
Evolution of plant anatomy
Arbuscular mycorrhizal symbiosis
There is no evidence that early land plants of the Silurian and early Devonian had roots, although fossil evidence of rhizoids occurs for several species, such as Horneophyton. The earliest land plants did not have vascular systems for transport of water and nutrients either. Aglaophyton, a rootless vascular plant known from Devonian fossils in the Rhynie chert was the first land plant discovered to have had a symbiotic relationship with fungi which formed arbuscular mycorrhizas, literally "tree-like fungal roots", in a well-defined cylinder of cells (ring in cross section) in the cortex of its stems. The fungi fed on the plant's sugars, in exchange for nutrients generated or extracted from the soil (especially phosphate), to which the plant would otherwise have had no access. Like other rootless land plants of the Silurian and early Devonian Aglaophyton may have relied on arbuscular mycorrhizal fungi for acquisition of water and nutrients from the soil.
The fungi were of the phylum Glomeromycota, a group that probably first appeared 1 billion years ago and still forms arbuscular mycorrhizal associations today with all major land plant groups from bryophytes to pteridophytes, gymnosperms and angiosperms and with more than 80% of vascular plants.
Evidence from DNA sequence analysis indicates that the arbuscular mycorrhizal mutualism arose in the common ancestor of these land plant groups during their transition to land and it may even have been the critical step that enabled them to colonise the land. Appearing as they did before these plants had evolved roots, mycorrhizal fungi would have assisted plants in the acquisition of water and mineral nutrients such as phosphorus, in exchange for organic compounds which they could not synthesize themselves. Such fungi increase the productivity even of simple plants such as liverworts.
Cuticle, stomata and intercellular spaces
To photosynthesise, plants must absorb from the atmosphere. However, making the tissues available for to enter allows water to evaporate, so this comes at a price. Water is lost much faster than is absorbed, so plants need to replace it. Early land plants transported water apoplastically, within the porous walls of their cells. Later, they evolved three anatomical features that provided the ability to control the inevitable water loss that accompanied acquisition. First, a waterproof outer covering or cuticle evolved that reduced water loss. Secondly, variable apertures, the stomata that could open and close to regulate the amount of water lost by evaporation during uptake and thirdly intercellular space between photosynthetic parenchyma cells that allowed improved internal distribution of the to the chloroplasts. This three-part system provided improved homoiohydry, the regulation of water content of the tissues, providing a particular advantage when water supply is not constant. The high concentrations of the Silurian and early Devonian, when plants were first colonising land, meant that they used water relatively efficiently. As was withdrawn from the atmosphere by plants, more water was lost in its capture, and more elegant water acquisition and transport mechanisms evolved. Plants growing upwards into the air needed a system for transporting water from the soil to all the different parts of the above-soil plant, especially to photosynthesising parts. By the end of the Carboniferous, when concentrations had been reduced to something approaching that of today, around 17 times more water was lost per unit of uptake. However, even in the "easy" early days, water was always at a premium, and had to be transported to parts of the plant from the wet soil to avoid desiccation.
Water can be wicked by capillary action along a fabric with small spaces. In narrow columns of water, such as those within the plant cell walls or in tracheids, when molecules evaporate from one end, they pull the molecules behind them along the channels. Therefore, evaporation alone provides the driving force for water transport in plants. However, without specialized transport vessels, this cohesion-tension mechanism can cause negative pressures sufficient to collapse water conducting cells, limiting the transport water to no more than a few cm, and therefore limiting the size of the earliest plants.
Xylem
To be free from the constraints of small size and constant moisture that the parenchymatic transport system inflicted, plants needed a more efficient water transport system. As plants grew upwards, specialised water transport vascular tissues evolved, first in the form of simple hydroids of the type found in the setae of moss sporophytes. These simple elongated cells were dead and water-filled at maturity, providing a channel for water transport, but their thin, unreinforced walls would collapse under modest water tension, limiting the plant height. Xylem tracheids, wider cells with lignin-reinforced cell walls that were more resistant to collapse under the tension caused by water stress, occur in more than one plant group by mid-Silurian, and may have a single evolutionary origin, possibly within the hornworts, uniting all tracheophytes. Alternatively, they may have evolved more than once. Much later, in the Cretaceous, tracheids were followed by vessels in flowering plants. As water transport mechanisms and waterproof cuticles evolved, plants could survive without being continually covered by a film of water. This transition from poikilohydry to homoiohydry opened up new potential for colonisation.
The early Devonian pretracheophytes Aglaophyton and Horneophyton have unreinforced water transport tubes with wall structures very similar to moss hydroids, but they grew alongside several species of tracheophytes, such as Rhynia gwynne-vaughanii that had xylem tracheids that were well reinforced by bands of lignin. The earliest macrofossils known to have xylem tracheids are small, mid-Silurian plants of the genus Cooksonia. However, thickened bands on the walls of isolated tube fragments are apparent from the early Silurian onwards.
Plants continued to innovate ways of reducing the resistance to flow within their cells, progressively increasing the efficiency of their water transport and to increase the resistance of the tracheids to collapse under tension. During the early Devonian, maximum tracheid diameter increased with time, but may have plateaued in the zosterophylls by mid-Devonian. Overall transport rate also depends on the overall cross-sectional area of the xylem bundle itself, and some mid-Devonian plants, such as the Trimerophytes, had much larger steles than their early ancestors. While wider tracheids provided higher rates of water transport, they increased the risk of cavitation, the formation of air bubbles resulting from the breakage of the water column under tension. Small pits in tracheid walls allow water to by-pass a defective tracheid while preventing air bubbles from passing through but at the cost of restricted flow rates. By the Carboniferous, Gymnosperms had developed bordered pits, valve-like structures that allow high-conductivity pits to seal when one side of a tracheid is depressurized.
Tracheids have non-perforated end walls with pits, which impose a great deal of resistance on water flow, but may have had the advantage of isolating air embolisms caused by cavitation or freezing. Vessels first evolved during the dry, low periods of the Late Permian, in the horsetails, ferns and Selaginellales independently, and later appeared in the mid Cretaceous in gnetophytes and angiosperms. Vessel members are open tubes with no end walls, and are arranged end to end to operate as if they were one continuous vessel. Vessels allowed the same cross-sectional area of wood to transport much more water than tracheids. This allowed plants to fill more of their stems with structural fibres and also opened a new niche to vines, which could transport water without being as thick as the tree they grew on. Despite these advantages, tracheid-based wood is a lot lighter, thus cheaper to make, as vessels need to be much more reinforced to avoid cavitation. Once plants had evolved this level of control over water evaporation and water transport, they were truly homoiohydric, able to extract water from their environment through root-like organs rather than relying on a film of surface moisture, enabling them to grow to much greater size but as a result of their increased independence from their surroundings, most vascular plants lost their ability to survive desiccation - a costly trait to lose. In early land plants, support was mainly provided by turgor pressure, particularly of the outer layer of cells known as the sterome tracheids, and not by the xylem, which was too small, too weak and in too central a position to provide much structural support. Plants with secondary xylem that had appeared by mid-Devonian, such as the Trimerophytes and Progymnosperms had much larger vascular cross sections producing strong woody tissue.
Endodermis
An endodermis may have evolved in the earliest plant roots during the Devonian, but the first fossil evidence for such a structure is Carboniferous. The endodermis in the roots surrounds the water transport tissue and regulates ion exchange between the groundwater and the tissues and prevents unwanted pathogens etc. from entering the water transport system. The endodermis can also provide an upwards pressure, forcing water out of the roots when transpiration is not enough of a driver.
Evolution of plant morphology
Leaves
Leaves are the primary photosynthetic organs of a modern plant. The origin of leaves was almost certainly triggered by falling concentrations of atmospheric during the Devonian period, increasing the efficiency with which carbon dioxide could be captured for photosynthesis.
Leaves evolved more than once. Based on their structure, they are classified into two types: microphylls, which lack complex venation and may have originated as spiny outgrowths known as enations, and megaphylls, which are large and have complex venation that may have arisen from the modification of groups of branches. It has been proposed that these structures arose independently. Megaphylls, according to Walter Zimmerman's telome theory, have evolved from plants that showed a three-dimensional branching architecture, through three transformations—overtopping, which led to the lateral position typical of leaves, planation, which involved formation of a planar architecture, webbing or fusion, which united the planar branches, thus leading to the formation of a proper leaf lamina. All three steps happened multiple times in the evolution of today's leaves.
It is widely believed that the telome theory is well supported by fossil evidence. However, Wolfgang Hagemann questioned it for morphological and ecological reasons and proposed an alternative theory. Whereas according to the telome theory the most primitive land plants have a three-dimensional branching system of radially symmetrical axes (telomes), according to Hagemann's alternative the opposite is proposed: the most primitive land plants that gave rise to vascular plants were flat, thalloid, leaf-like, without axes, somewhat like a liverwort or fern prothallus. Axes such as stems and roots evolved later as new organs. Rolf Sattler proposed an overarching process-oriented view that leaves some limited room for both the telome theory and Hagemann's alternative and in addition takes into consideration the whole continuum between dorsiventral (flat) and radial (cylindrical) structures that can be found in fossil and living land plants. This view is supported by research in molecular genetics. Thus, James (2009) concluded that "it is now widely accepted that... radiality [characteristic of axes such as stems] and dorsiventrality [characteristic of leaves] are but extremes of a continuous spectrum. In fact, it is simply the timing of the KNOX gene expression".
Before the evolution of leaves, plants had the photosynthetic apparatus on the stems, which they retain albeit leaves have largely assumed that job. Today's megaphyll leaves probably became commonplace some 360mya, about 40my after the simple leafless plants had colonized the land in the Early Devonian. This spread has been linked to the fall in the atmospheric carbon dioxide concentrations in the Late Paleozoic era associated with a rise in density of stomata on leaf surface. This would have resulted in greater transpiration rates and gas exchange, but especially at high concentrations, large leaves with fewer stomata would have heated to lethal temperatures in full sunlight. Increasing the stomatal density allowed for a better-cooled leaf, thus making its spread feasible, but increased uptake at the expense of decreased water use efficiency.
The rhyniophytes of the Rhynie chert consisted only of slender, unornamented axes. The early to middle Devonian trimerophytes may be considered leafy. This group of vascular plants are recognisable by their masses of terminal sporangia, which adorn the ends of axes which may bifurcate or trifurcate. Some organisms, such as Psilophyton, bore enations. These are small, spiny outgrowths of the stem, lacking their own vascular supply.
The zosterophylls were already important in the late Silurian, much earlier than any rhyniophytes of comparable complexity. This group, recognisable by their kidney-shaped sporangia which grew on short lateral branches close to the main axes, sometimes branched in a distinctive H-shape. Many zosterophylls bore enations (small tissue outgrowths on the surface with variable morphologies) on their axes but none of these had a vascular trace. The first evidence of vascularised enations occurs in a fossil clubmoss known as Baragwanathia that had already appeared in the fossil record in the Late Silurian. In this organism, these leaf traces continue into the leaf to form their mid-vein. One theory, the "enation theory", holds that the microphyllous leaves of clubmosses developed by outgrowths of the protostele connecting with existing enations The leaves of the Rhynie genus Asteroxylon, which was preserved in the Rhynie chert almost 20 million years later than Baragwanathia, had a primitive vascular supply – in the form of leaf traces departing from the central protostele towards each individual "leaf". Asteroxylon and Baragwanathia are widely regarded as primitive lycopods, a group still extant today, represented by the quillworts, the spikemosses and the club mosses. Lycopods bear distinctive microphylls, defined as leaves with a single vascular trace. Microphylls could grow to some size, those of Lepidodendrales reaching over a meter in length, but almost all just bear the one vascular bundle. An exception is the rare branching in some Selaginella species.
The more familiar leaves, megaphylls, are thought to have originated four times independently: in the ferns, horsetails, progymnosperms and seed plants. They appear to have originated by modifying dichotomising branches, which first overlapped (or "overtopped") one another, became flattened or planated and eventually developed "webbing" and evolved gradually into more leaf-like structures. Megaphylls, by Zimmerman's telome theory, are composed of a group of webbed branches and hence the "leaf gap" left where the leaf's vascular bundle leaves that of the main branch resembles two axes splitting. In each of the four groups to evolve megaphylls, their leaves first evolved during the Late Devonian to Early Carboniferous, diversifying rapidly until the designs settled down in the mid Carboniferous.
The cessation of further diversification can be attributed to developmental constraints, raising the question of why it took so long for leaves to evolve in the first place. Plants had been on land for at least 50 million years before megaphylls became significant. However, small, rare mesophylls are known from the early Devonian genus Eophyllophyton – so development could not have been a barrier to their appearance. The best explanation so far is that atmospheric was declining rapidly during this time – falling by around 90% during the Devonian. This required an increase in stomatal density by 100 times to maintain the rate of photosynthesis. When stomata open to allow water to evaporate from leaves it has a cooling effect, resulting from the loss of latent heat of evaporation. It appears that the low stomatal density in the early Devonian meant that evaporation and evaporative cooling were limited, and that leaves would have overheated if they grew to any size. The stomatal density could not increase, as the primitive steles and limited root systems would not be able to supply water quickly enough to match the rate of transpiration. Clearly, leaves are not always beneficial, as illustrated by the frequent occurrence of secondary loss of leaves, exemplified by cacti and the "whisk fern" Psilotum.
Secondary evolution can disguise the true evolutionary origin of some leaves. Some genera of ferns display complex leaves which are attached to the pseudostele by an outgrowth of the vascular bundle, leaving no leaf gap.
Deciduous trees deal with another disadvantage to having leaves. The popular belief that plants shed their leaves when the days get too short is misguided; evergreens prospered in the Arctic Circle during the most recent greenhouse earth. The generally accepted reason for shedding leaves during winter is to cope with the weather – the force of wind and weight of snow are much more comfortably weathered without leaves to increase surface area. Seasonal leaf loss has evolved independently several times and is exhibited in the ginkgoales, some pinophyta and certain angiosperms. Leaf loss may also have arisen as a response to pressure from insects; it may have been less costly to lose leaves entirely during the winter or dry season than to continue investing resources in their repair.
Roots
The evolution of roots had consequences on a global scale. By disturbing the soil and promoting its acidification (by taking up nutrients such as nitrate and phosphate), they enabled it to weather more deeply, injecting carbon compounds deeper into soils with huge implications for climate. These effects may have been so profound they led to a mass extinction.
While there are traces of root-like impressions in fossil soils in the Late Silurian, body fossils show the earliest plants to be devoid of roots. Many had prostrate branches that sprawled along the ground, with upright axes or thalli dotted here and there, and some even had non-photosynthetic subterranean branches which lacked stomata. Roots have a root cap, unlike specialised branches. So while Siluro-Devonian plants such as Rhynia and Horneophyton possessed the physiological equivalent of roots, roots – defined as organs differentiated from stems – did not arrive until later. Unfortunately, roots are rarely preserved in the fossil record.
Rhizoids – small structures performing the same role as roots, usually a cell in diameter – probably evolved very early, perhaps even before plants colonised the land; they are recognised in the Characeae, an algal sister group to land plants. That said, rhizoids probably evolved more than once; the rhizines of lichens, for example, perform a similar role. Even some animals (Lamellibrachia) have root-like structures. Rhizoids are clearly visible in the Rhynie chert fossils, and were present in most of the earliest vascular plants, and on this basis seem to have presaged true plant roots.
More advanced structures are common in the Rhynie chert, and many other fossils of comparable early Devonian age bear structures that look like, and acted like, roots. The rhyniophytes bore fine rhizoids, and the trimerophytes and herbaceous lycopods of the chert bore root-like structure penetrating a few centimetres into the soil. However, none of these fossils display all the features borne by modern roots, with the exception of Asteroxylon, which has recently been recognized as bearing roots that evolved independently from those of extant vascular plants. Roots and root-like structures became increasingly common and deeper penetrating during the Devonian, with lycopod trees forming roots around 20 cm long during the Eifelian and Givetian. These were joined by progymnosperms, which rooted up to about a metre deep, during the ensuing Frasnian stage. True gymnosperms and zygopterid ferns also formed shallow rooting systems during the Famennian.
The rhizophores of the lycopods provide a slightly different approach to rooting. They were equivalent to stems, with organs equivalent to leaves performing the role of rootlets. A similar construction is observed in the extant lycopod Isoetes, and this appears to be evidence that roots evolved independently at least twice, in the lycophytes and other plants, a proposition supported by studies showing that roots are initiated and their growth promoted by different mechanisms in lycophytes and euphyllophytes.
Early rooted plants are little more advanced than their Silurian forebears, without a dedicated root system; however, the flat-lying axes can be clearly seen to have growths similar to the rhizoids of bryophytes today.
By the Middle to Late Devonian, most groups of plants had independently developed a rooting system of some nature. As roots became larger, they could support larger trees, and the soil was weathered to a greater depth. This deeper weathering had effects not only on the aforementioned drawdown of , but also opened up new habitats for colonisation by fungi and animals.
The narrowest roots of modern plants are a mere 40 μm in diameter, and could not physically transport water if they were any narrower. The earliest fossil roots recovered, by contrast, narrowed from 3 mm to under 700 μm in diameter; of course, taphonomy is the ultimate control of what thickness can be seen.
Tree form
The early Devonian landscape was devoid of vegetation taller than waist height. Greater height provided a competitive advantage in the harvesting of sunlight for photosynthesis, overshadowing of competitors and in spore distribution, as spores (and later, seeds) could be blown for greater distances if they started higher. An effective vascular system was required in order to achieve greater heights. To attain arborescence, plants had to develop woody tissue that provided both support and water transport, and thus needed to evolve the capacity for secondary growth. The stele of plants undergoing secondary growth is surrounded by a vascular cambium, a ring of meristematic cells which produces more xylem on the inside and phloem on the outside. Since xylem cells comprise dead, lignified tissue, subsequent rings of xylem are added to those already present, forming wood. Fossils of plants from the early Devonian show that a simple form of wood first appeared at least 400 million years ago, at a time when all land plants were small and herbaceous. Because wood evolved long before shrubs and trees, it is likely that its original purpose was for water transport, and that it was only used for mechanical support later.
The first plants to develop secondary growth and a woody habit, were apparently the ferns, and as early as the Middle Devonian one species, Wattieza, had already reached heights of 8 m and a tree-like habit.
Other clades did not take long to develop a tree-like stature. The Late Devonian Archaeopteris, a precursor to gymnosperms which evolved from the trimerophytes, reached 30 m in height. The progymnosperms were the first plants to develop true wood, grown from a bifacial cambium. The first appearance of one of them, Rellimia, was in the Middle Devonian. True wood is only thought to have evolved once, giving rise to the concept of a "lignophyte" clade.
Archaeopteris forests were soon supplemented by arborescent lycopods, in the form of Lepidodendrales, which exceeded 50m in height and 2m across at the base. These arborescent lycopods rose to dominate Late Devonian and Carboniferous forests that gave rise to coal deposits. Lepidodendrales differ from modern trees in exhibiting determinate growth: after building up a reserve of nutrients at a lower height, the plants would "bolt" as a single trunk to a genetically determined height, branch at that level, spread their spores and die. They consisted of "cheap" wood to allow their rapid growth, with at least half of their stems comprising a pith-filled cavity. Their wood was also generated by a unifacial vascular cambium – it did not produce new phloem, meaning that the trunks could not grow wider over time.
The horsetail Calamites appeared in the Carboniferous. Unlike the modern horsetail Equisetum, Calamites had a unifacial vascular cambium, allowing them to develop wood and grow to heights in excess of 10 m and to branch repeatedly.
While the form of early trees was similar to that of today's, the Spermatophytes or seed plants, the group that contain all modern trees, had yet to evolve. The dominant tree groups today are all seed plants, the gymnosperms, which include the coniferous trees, and the angiosperms, which contain all fruiting and flowering trees. No free-sporing trees like Archaeopteris exist in the extant flora. It was long thought that the angiosperms arose from within the gymnosperms, but recent molecular evidence suggests that their living representatives form two distinct groups. The molecular data has yet to be fully reconciled with morphological data, but it is becoming accepted that the morphological support for paraphyly is not especially strong.
This would lead to the conclusion that both groups arose from within the pteridosperms, probably as early as the Permian.
The angiosperms and their ancestors played a very small role until they diversified during the Cretaceous. They started out as small, damp-loving organisms in the understorey, and have been diversifying ever since the Cretaceous, to become the dominant member of non-boreal forests today.
Seeds
Early land plants reproduced in the fashion of ferns: spores germinated into small gametophytes, which produced eggs and/or sperm. These sperm would swim across moist soils to find the female organs (archegonia) on the same or another gametophyte, where they would fuse with an egg to produce an embryo, which would germinate into a sporophyte.
Heterosporic plants, as their name suggests, bear spores of two sizes – microspores and megaspores. These would germinate to form microgametophytes and megagametophytes, respectively. This system paved the way for ovules and seeds: taken to the extreme, the megasporangia could bear only a single megaspore tetrad, and to complete the transition to true ovules, three of the megaspores in the original tetrad could be aborted, leaving one megaspore per megasporangium.
The transition to ovules continued with this megaspore being "boxed in" to its sporangium while it germinated. Then, the megagametophyte was contained within a waterproof integument, which enclosed the seed. The pollen grain, which contained a microgametophyte germinated from a microspore , was employed for dispersal of the male gamete, only releasing its desiccation-prone flagellate sperm when it reached a receptive megagametophyte.
Lycopods and sphenopsids got a fair way down the path to the seed habit without ever crossing the threshold. Fossil lycopod megaspores reaching 1 cm in diameter, and surrounded by vegetative tissue, are known (Lepidocarpon, Achlamydocarpon);– these even germinated into a megagametophyte in situ. However, they fell short of being ovules, since the nucellus, an inner spore-covering layer, does not completely enclose the spore. A very small slit (micropyle) remains, meaning that the megasporangium is still exposed to the atmosphere. This has two consequences – firstly, it means it is not fully resistant to desiccation, and secondly, sperm do not have to "burrow" to access the archegonia of the megaspore.
A Middle Devonian precursor to seed plants from Belgium has been identified predating the earliest seed plants by about 20 million years. Runcaria, small and radially symmetrical, is an integumented megasporangium surrounded by a cupule. The megasporangium bears an unopened distal extension protruding above the multilobed integument. It is suspected that the extension was involved in anemophilous pollination. Runcaria sheds new light on the sequence of character acquisition leading to the seed. Runcaria has all of the qualities of seed plants except for a solid seed coat and a system to guide the pollen to the ovule.
The first spermatophytes (literally: "seed plants") – that is, the first plants to bear true seeds – are called pteridosperms: literally, "seed ferns", so called because their foliage consisted of fern-like fronds, although they were not closely related to ferns. The oldest fossil evidence of seed plants is of Late Devonian age, and they appear to have evolved out of an earlier group known as the progymnosperms. These early seed plants ranged from trees to small, rambling shrubs; like most early progymnosperms, they were woody plants with fern-like foliage. They all bore ovules, but no cones, fruit or similar. While it is difficult to track the early evolution of seeds, the lineage of the seed ferns may be traced from the simple trimerophytes through homosporous Aneurophytes.
The seed plants underwent their first major evolutionary radiation in the Famennian era.
This seed model is shared by basically all gymnosperms (literally: "naked seeds"), most of which encase their seeds in a woody cone or fleshy aril (the yew, for example), but none of which fully enclose their seeds. The angiosperms ("vessel seeds") are the only group to fully enclose the seed, in a carpel.
Fully enclosed seeds opened up a new pathway for plants to follow: that of seed dormancy. The embryo, completely isolated from the external atmosphere and hence protected from desiccation, could survive some years of drought before germinating.
Gymnosperm seeds from the Late Carboniferous have been found to contain embryos, suggesting a lengthy gap between fertilisation and germination. This period is associated with the entry into a greenhouse earth period, with an associated increase in aridity. This suggests that dormancy arose as a response to drier climatic conditions, where it became advantageous to wait for a moist period before germinating. This evolutionary breakthrough appears to have opened a floodgate: previously inhospitable areas, such as dry mountain slopes, could now be tolerated, and were soon covered by trees.
Seeds offered further advantages to their bearers: they increased the success rate of fertilised gametophytes, and because a nutrient store could be "packaged" in with the embryo, the seeds could germinate rapidly in inhospitable environments, reaching a size where it could fend for itself more quickly. For example, without an endosperm, seedlings growing in arid environments would not have the reserves to grow roots deep enough to reach the water table before they expired from dehydration. Likewise, seeds germinating in a gloomy understory require an additional reserve of energy to quickly grow high enough to capture sufficient light for self-sustenance.
A combination of these advantages gave seed plants the ecological edge over the previously dominant genus Archaeopteris, thus increasing the biodiversity of early forests.
Despite these advantages, it is common for fertilized ovules to fail to mature as seeds. Also during seed dormancy (often associated with unpredictable and stressful conditions) DNA damage accumulates. Thus DNA damage appears to be a basic problem for survival of seed plants, just as DNA damage is a major problem for life in general.
Flowers
Flowers are modified leaves possessed only by the angiosperms, which are relatively late to appear in the fossil record. The group originated and diversified during the Early Cretaceous and became ecologically significant thereafter. Flower-like structures first appear in the fossil records some ~130 mya, in the Cretaceous. However, in 2018, scientists reported the finding of a fossil flower from about 180 million years ago, 50 million years earlier than previously thought. The interpretation has been however highly disputed.
Colorful and/or pungent structures surround the cones of plants such as cycads and Gnetales, making a strict definition of the term "flower" elusive.
The main function of a flower is reproduction, which, before the evolution of the flower and angiosperms, was the job of microsporophylls and megasporophylls. A flower can be considered a powerful evolutionary innovation, because its presence allowed the plant world to access new means and mechanisms for reproduction.
The flowering plants have long been assumed to have evolved from within the gymnosperms; according to the traditional morphological view, they are closely allied to the Gnetales. However, as noted above, recent molecular evidence is at odds with this hypothesis, and further suggests that Gnetales are more closely related to some gymnosperm groups than angiosperms, and that extant gymnosperms form a distinct clade to the angiosperms, the two clades diverging some .
The relationship of stem groups to the angiosperms is important in determining the evolution of flowers. Stem groups provide an insight into the state of earlier "forks" on the path to the current state. Convergence increases the risk of misidentifying stem groups. Since the protection of the megagametophyte is evolutionarily desirable, probably many separate groups evolved protective encasements independently. In flowers, this protection takes the form of a carpel, evolved from a leaf and recruited into a protective role, shielding the ovules. These ovules are further protected by a double-walled integument.
Penetration of these protective layers needs something more than a free-floating microgametophyte. Angiosperms have pollen grains comprising just three cells. One cell is responsible for drilling down through the integuments, and creating a conduit for the two sperm cells to flow down. The megagametophyte has just seven cells; of these, one fuses with a sperm cell, forming the nucleus of the egg itself, and another joins with the other sperm, and dedicates itself to forming a nutrient-rich endosperm. The other cells take auxiliary roles. This process of "double fertilisation" is unique and common to all angiosperms.
In the fossil record, there are three intriguing groups which bore flower-like structures. The first is the Permian pteridosperm Glossopteris, which already bore recurved leaves resembling carpels. The Mesozoic Caytonia is more flower-like still, with enclosed ovules – but only a single integument. Further, details of their pollen and stamens set them apart from true flowering plants.
The Bennettitales bore remarkably flower-like organs, protected by whorls of bracts which may have played a similar role to the petals and sepals of true flowers; however, these flower-like structures evolved independently, as the Bennettitales are more closely related to cycads and ginkgos than to the angiosperms.
However, no true flowers are found in any groups save those extant today. Most morphological and molecular analyses place Amborella, the nymphaeales and Austrobaileyaceae in a basal clade called "ANA". This clade appear to have diverged in the early Cretaceous, around – around the same time as the earliest fossil angiosperm, and just after the first angiosperm-like pollen, 136 million years ago. The magnoliids diverged soon after, and a rapid radiation had produced eudicots and monocots by . By the end of the Cretaceous , over 50% of today's angiosperm orders had evolved, and the clade accounted for 70% of global species. It was around this time that flowering trees became dominant over conifers.
The features of the basal "ANA" groups suggest that angiosperms originated in dark, damp, frequently disturbed areas. It appears that the angiosperms remained constrained to such habitats throughout the Cretaceous – occupying the niche of small herbs early in the successional series. This may have restricted their initial significance, but given them the flexibility that accounted for the rapidity of their later diversifications in other habitats.
Some propose that the Angiosperms arose from an unknown Seed Fern, Pteridophyte, and view Cycads as living Seed Ferns with both Seed-Bearing and sterile leaves (Cycas revoluta)
In August 2017, scientists presented a detailed description and 3D reconstruction of possibly the first flower that lived about 140 million years ago.
Origins of the flower
The family Amborellaceae is regarded as being the sister clade to all other living flowering plants. A draft genome of Amborella trichopoda was published in December, 2013. By comparing its genome with those of all other living flowering plants, it will be possible to work out the most likely characteristics of the ancestor of A. trichopoda and all other flowering plants, i.e. the ancestral flowering plant.
It seems that on the level of the organ, the leaf may be the ancestor of the flower, or at least some floral organs. When some crucial genes involved in flower development are mutated, clusters of leaf-like structures arise in place of flowers. Thus, sometime in history, the developmental program leading to formation of a leaf must have been altered to generate a flower. There probably also exists an overall robust framework within which the floral diversity has been generated. An example of that is a gene called LEAFY (LFY), which is involved in flower development in Arabidopsis thaliana. The homologs of this gene are found in angiosperms as diverse as tomato, snapdragon, pea, maize and even gymnosperms. Expression of Arabidopsis thaliana LFY in distant plants like poplar and citrus also results in flower-production in these plants. The LFY gene regulates the expression of some genes belonging to the MADS-box family. These genes, in turn, act as direct controllers of flower development.
Adaptive function of flowers
Flowers likely emerged during plant evolution as an adaptation to facilitate cross-fertilization (outcrossing), a process that leads to the masking of recessive deleterious mutations in progeny genomes. This masking effect of expression of deleterious mutations is referred to as genetic complementation. This beneficial masking effect of cross-fertilization is also considered to be the basis of hybrid vigor or heterosis in progeny. Once flowers have become established in a lineage based on their adaptive function of promoting cross-fertilization, subsequent switching to inbreeding ordinarily then becomes disadvantageous, mainly because it permits expression of the previously masked deleterious recessive mutations, i.e. inbreeding depression. In addition, meiosis, the process by which seed progeny are produced in flowering plants, provides a direct mechanism for repairing DNA through genetic recombination. Thus, in flowering plants, the two fundamental processes of sexual reproduction are cross-fertilization (outcrossing) and meiosis and these two processes appear to be maintained respectively by the advantages of genetic complementation and recombinational repair of DNA.
Evolution of the MADS-box family
The members of the MADS-box family of transcription factors play a very important and evolutionarily conserved role in flower development. According to the ABC Model of flower development, three zones — A, B and C — are generated within the developing flower primordium, by the action of some transcription factors, that are members of the MADS-box family. Among these, the functions of the B and C domain genes have been evolutionarily more conserved than the A domain gene. Many of these genes have arisen through gene duplications of ancestral members of this family. Quite a few of them show redundant functions.
The evolution of the MADS-box family has been extensively studied. These genes are present even in pteridophytes, but the spread and diversity is many times higher in angiosperms. There appears to be quite a bit of pattern into how this family has evolved. Consider the evolution of the C-region gene AGAMOUS (AG). It is expressed in today's flowers in the stamens, and the carpel, which are reproductive organs. Its ancestor in gymnosperms also has the same expression pattern. Here, it is expressed in the strobili, an organ that produces pollen or ovules. Similarly, the B-genes' (AP3 and PI) ancestors are expressed only in the male organs in gymnosperms. Their descendants in the modern angiosperms also are expressed only in the stamens, the male reproductive organ. Thus, the same, then-existing components were used by the plants in a novel manner to generate the first flower. This is a recurring pattern in evolution.
Factors influencing floral diversity
There is enormous variation in floral structure in plants, typically due to changes in the MADS-box genes and their expression pattern. For example, grasses possess unique floral structures. The carpels and stamens are surrounded by scale-like lodicules and two bracts, the lemma and the palea, but genetic evidence and morphology suggest that lodicules are homologous to eudicot petals. The palea and lemma may be homologous to sepals in other groups, or may be unique grass structures. Another example is that of Linaria vulgaris, which has two kinds of flower symmetries-radial and bilateral. These symmetries are due to epigenetic changes in just one gene called CYCLOIDEA.
Arabidopsis thaliana has a gene called AGAMOUS that plays an important role in defining how many petals and sepals and other organs are generated. Mutations in this gene give rise to the floral meristem obtaining an indeterminate fate, and proliferation of floral organs in double-flowered forms of roses, carnations and morning glory. These phenotypes have been selected by horticulturists for their increased number of petals. Several studies on diverse plants like petunia, tomato, Impatiens, maize, etc. have suggested that the enormous diversity of flowers is a result of small changes in genes controlling their development.
The Floral Genome Project confirmed that the ABC Model of flower development is not conserved across all angiosperms. Sometimes expression domains change, as in the case of many monocots, and also in some basal angiosperms like Amborella. Different models of flower development like the Fading boundaries model, or the Overlapping-boundaries model which propose non-rigid domains of expression, may explain these architectures. There is a possibility that from the basal to the modern angiosperms, the domains of floral architecture have become more and more fixed through evolution.
Flowering time
Another floral feature that has been a subject of natural selection is flowering time. Some plants flower early in their life cycle, others require a period of vernalization before flowering. This outcome is based on factors like temperature, light intensity, presence of pollinators and other environmental signals: genes like CONSTANS (CO), Flowering Locus C (FLC) and FRIGIDA regulate integration of environmental signals into the pathway for flower development. Variations in these loci have been associated with flowering time variations between plants. For example, Arabidopsis thaliana ecotypes that grow in the cold, temperate regions require prolonged vernalization before they flower, while the tropical varieties, and the most common lab strains, don't. This variation is due to mutations in the FLC and FRIGIDA genes, rendering them non-functional.
Many of the genes involved in this process are conserved across all the plants studied. Sometimes though, despite genetic conservation, the mechanism of action turns out to be different. For example, rice is a short-day plant, while Arabidopsis thaliana is a long-day plant. Both plants have the proteins CO and FLOWERING LOCUS T (FT), but, in Arabidopsis thaliana, CO enhances FT production, while in rice, the CO homolog represses FT production, resulting in completely opposite downstream effects.
Theories of flower evolution
The Anthophyte theory was based on the observation that a gymnospermic group Gnetales has a flower-like ovule. It has partially developed vessels as found in the angiosperms, and the megasporangium is covered by three envelopes, like the ovary structure of angiosperm flowers. However, many other lines of evidence show that Gnetales is not related to angiosperms.
The Mostly Male theory has a more genetic basis. Proponents of this theory point out that the gymnosperms have two very similar copies of the gene LFY, while angiosperms just have one. Molecular clock analysis has shown that the other LFY paralog was lost in angiosperms around the same time as flower fossils become abundant, suggesting that this event might have led to floral evolution. According to this theory, loss of one of the LFY paralog led to flowers that were more male, with the ovules being expressed ectopically. These ovules initially performed the function of attracting pollinators, but sometime later, may have been integrated into the core flower.
Mechanisms and players in evolution of plant morphology
While environmental factors are significantly responsible for evolutionary change, they act merely as agents for natural selection. Change is inherently brought about via phenomena at the genetic level: mutations, chromosomal rearrangements, and epigenetic changes. While the general types of mutations hold true across the living world, in plants, some other mechanisms have been implicated as highly significant.
Genome doubling is a relatively common occurrence in plant evolution and results in polyploidy, which is consequently a common feature in plants. It is estimated that at least half (and probably all) plants have seen genome doubling in their history. Genome doubling entails gene duplication, thus generating functional redundancy in most genes. The duplicated genes may attain new function, either by changes in expression pattern or changes in activity. Polyploidy and gene duplication are believed to be among the most powerful forces in evolution of plant form; though it is not known why genome doubling is such a frequent process in plants. One probable reason is the production of large amounts of secondary metabolites in plant cells. Some of them might interfere in the normal process of chromosomal segregation, causing genome duplication.
In recent times, plants have been shown to possess significant microRNA families, which are conserved across many plant lineages. In comparison to animals, while the number of plant miRNA families are lesser than animals, the size of each family is much larger. The miRNA genes are also much more spread out in the genome than those in animals, where they are more clustered. It has been proposed that these miRNA families have expanded by duplications of chromosomal regions. Many miRNA genes involved in regulation of plant development have been found to be quite conserved between plants studied.
Domestication of plants like maize, rice, barley, wheat etc. has also been a significant driving force in their evolution. Research concerning the origin of maize has found that it is a domesticated derivative of a wild plant from Mexico called teosinte. Teosinte belongs to the genus Zea, just as maize, but bears very small inflorescence, 5–10 hard cobs and a highly branched and spread out stem.
Crosses between a particular teosinte variety and maize yields fertile offspring that are intermediate in phenotype between maize and teosinte. QTL analysis has also revealed some loci that, when mutated in maize, yield a teosinte-like stem or teosinte-like cobs. Molecular clock analysis of these genes estimates their origins to some 9,000 years ago, well in accordance with other records of maize domestication. It is believed that a small group of farmers must have selected some maize-like natural mutant of teosinte some 9,000 years ago in Mexico, and subjected it to continuous selection to yield the familiar maize plant of today.
The edible cauliflower is a domesticated version of the wild plant Brassica oleracea, which does not possess the dense undifferentiated inflorescence, called the curd, that cauliflower possesses.
Cauliflower possesses a single mutation in a gene called CAL, controlling meristem differentiation into inflorescence. This causes the cells at the floral meristem to gain an undifferentiated identity and, instead of growing into a flower, they grow into a dense mass of inflorescence meristem cells in arrested development. This mutation has been selected through domestication since at least the time of the Greek empire.
Evolution of photosynthetic pathways
The C4 metabolic pathway is a valuable recent evolutionary innovation in plants, involving a complex set of adaptive changes to physiology and gene expression patterns.
Photosynthesis is a complex chemical pathway facilitated by a range of enzymes and co-enzymes. The enzyme RuBisCO is responsible for "fixing" – that is, it attaches it to a carbon-based molecule to form a sugar, which can be used by the plant, releasing an oxygen molecule. However, the enzyme is notoriously inefficient, and, as ambient temperature rises, will increasingly fix oxygen instead of in a process called photorespiration. This is energetically costly as the plant has to use energy to turn the products of photorespiration back into a form that can react with .
Concentrating carbon
Broadly, the two main ways to concentrate carbon dioxide in plants are 1) biochemical concentrating mechanisms (CCM) and 2) biophysical concentrating mechanisms. Biochemical CCMs such as C4 and CAM photosynthesis concentrate by using an enzyme, phosphoenolpyruvate carboxylase, to bind inorganic carbon to an intermediate four carbon sugar, which can then be converted back to RuBP and for subsequent fixation by Rubisco. Biophysical CCMs like carboxysomes and pyrenoids concentrate in a particular locus through the coordination of carbonic anhydrases and anion channels.
C4 plants evolved carbon concentrating mechanisms that work by increasing the concentration of around RuBisCO, and excluding oxygen, thereby increasing the efficiency of photosynthesis by decreasing photorespiration. The process of concentrating around RuBisCO requires more energy than allowing gases to diffuse, but under certain conditions – i.e. warm temperatures (>25 °C), low concentrations, or high oxygen concentrations – pays off in terms of the decreased loss of sugars through photorespiration.
One type of C4 metabolism employs a so-called Kranz anatomy. This transports through an outer mesophyll layer, via a range of organic molecules, to the central bundle sheath cells, where the is released. In this way, is concentrated near the site of RuBisCO operation. Because RuBisCO is operating in an environment with much more than it otherwise would be, it performs more efficiently.
A second mechanism, CAM photosynthesis, temporally separates photosynthesis from the action of RuBisCO. RuBisCO only operates during the day, when stomata are sealed and is provided by the breakdown of the chemical malate. More is then harvested from the atmosphere when stomata open, during the cool, moist nights, reducing water loss.
The third mechanism present in plants, pyrenoid-based CCMs, is found only in the hornwort lineage . In this mechanism, RuBisCO is concentrated in the pyrenoid, a membraneless compartment, by importing inorganic carbon in the form of bicarbonate . This import is thought to be dependent on the coordination of carbonic anhydrases and anion channels, and takes advantage of the native pH differences between the cytosol, chloroplast stroma, and thylakoid lumen.
Evolutionary record
These two pathways, with the same effect on RuBisCO, evolved a number of times independently – indeed, C4 alone arose 62 times in 18 different plant families. A number of 'pre-adaptations' seem to have paved the way for , leading to its clustering in certain clades: it has most frequently been innovated in plants that already had features such as extensive vascular bundle sheath tissue. Many potential evolutionary pathways resulting in the phenotype are possible and have been characterised using Bayesian inference, confirming that non-photosynthetic adaptations often provide evolutionary stepping stones for the further evolution of .
The C4 construction is used by a subset of grasses, while CAM is employed by many succulents and cacti. The C4 trait appears to have emerged during the Oligocene, around ; however, they did not become ecologically significant until the Miocene, . Remarkably, some charcoalified fossils preserve tissue organised into the Kranz anatomy, with intact bundle sheath cells, allowing the presence C4 metabolism to be identified. Isotopic markers are used to deduce their distribution and significance.
C3 plants preferentially use the lighter of two isotopes of carbon in the atmosphere, 12C, which is more readily involved in the chemical pathways involved in its fixation. Because C4 metabolism involves a further chemical step, this effect is accentuated. Plant material can be analysed to deduce the ratio of the heavier 13C to 12C. This ratio is denoted . C3 plants are on average around 14‰ (parts per thousand) lighter than the atmospheric ratio, while C4 plants are about 28‰ lighter. The of CAM plants depends on the percentage of carbon fixed at night relative to what is fixed in the day, being closer to C3 plants if they fix most carbon in the day and closer to C4 plants if they fix all their carbon at night.
Original fossil material in sufficient quantity to analyse the grass itself is scarce, but horses provide a good proxy. They were globally widespread in the period of interest, and browsed almost exclusively on grasses. There's an old phrase in isotope paleontology, "you are what you eat (plus a little bit)" – this refers to the fact that organisms reflect the isotopic composition of whatever they eat, plus a small adjustment factor. There is a good record of horse teeth throughout the globe, and their record shows a sharp negative inflection around , during the Messinian that is interpreted as resulting from the rise of C4 plants on a global scale.
Advantage of C4
While C4 enhances the efficiency of RuBisCO, the concentration of carbon is highly energy intensive. This means that C4 plants only have an advantage over C3 organisms in certain conditions: namely, high temperatures and low rainfall. C4 plants also need high levels of sunlight to thrive. Models suggest that, without wildfires removing shade-casting trees and shrubs, there would be no space for C4 plants. But, wildfires have occurred for 400 million years. The Carboniferous (~) had notoriously high oxygen levels – almost enough to allow spontaneous combustion – and very low , but no C4 isotopic signature has been found. There also does not seem to be a sudden trigger for the Miocene rise.
During the Miocene, the atmosphere and climate were relatively stable. If anything, increased gradually from before settling down to concentrations similar to the Holocene. This suggests that it did not have a key role in invoking C4 evolution. Grasses themselves (the group which would give rise to the most occurrences of C4) had probably been around for 60 million years or more, so had had plenty of time to evolve C4, which, in any case, is present in a diverse range of groups and thus evolved independently. There is a strong signal of climate change in South Asia; increasing aridity – hence increasing fire frequency and intensity – may have led to an increase in the importance of grasslands. However, this is difficult to reconcile with the North American record. It is possible that the signal is entirely biological, forced by the fire driven acceleration of grass evolution – which, both by increasing weathering and incorporating more carbon into sediments, reduced atmospheric levels. Finally, there is evidence that the onset of C4 from is a biased signal, which only holds true for North America, from where most samples originate; emerging evidence suggests that grasslands evolved to a dominant state at least 15Ma earlier in South America.
Evolution of transcriptional regulation
Transcription factors and transcriptional regulatory networks play key roles in plant development and stress responses, as well as their evolution. During plant landing, many novel transcription factor families emerged and are preferentially wired into the networks of multicellular development, reproduction, and organ development, contributing to more complex morphogenesis of land plants.
Evolution of secondary metabolism
Secondary metabolites are essentially low molecular weight compounds, sometimes having complex structures, that are not essential for the normal processes of growth, development, or reproduction. They function in processes as diverse as immunity, anti-herbivory, pollinator attraction, communication between plants, maintaining symbiotic associations with soil flora, or enhancing the rate of fertilization, and hence are significant from the evo-devo perspective. Secondary metabolites are structurally and functionally diverse, and it is estimated that hundreds of thousands of enzymes might be involved in the process of producing them, with about 15–25% of the genome coding for these enzymes, and every species having its unique arsenal of secondary metabolites. Many of these metabolites, such as salicylic acid are of medical significance to humans.
The purpose of producing so many secondary metabolites, with a significant proportion of the metabolome devoted to this activity is unclear. It is postulated that most of these chemicals help in generating immunity and, in consequence, the diversity of these metabolites is a result of a constant arms race between plants and their parasites. Some evidence supports this case. A central question involves the reproductive cost to maintaining such a large inventory of genes devoted to producing secondary metabolites. Various models have been suggested that probe into this aspect of the question, but a consensus on the extent of the cost has yet to be established; as it is still difficult to predict whether a plant with more secondary metabolites increases its survival or reproductive success compared to other plants in its vicinity.
Secondary metabolite production seems to have arisen quite early during evolution. In plants, they seem to have spread out using mechanisms including gene duplications or the evolution of novel genes. Furthermore, research has shown that diversity in some of these compounds may be positively selected for. Although the role of novel gene evolution in the evolution of secondary metabolism is clear, there are several examples where new metabolites have been formed by small changes in the reaction. For example, cyanogen glycosides have been proposed to have evolved multiple times in different plant lineages. There are several such instances of convergent evolution. For example, enzymes for synthesis of limonene – a terpene – are more similar between angiosperms and gymnosperms than to their own terpene synthesis enzymes. This suggests independent evolution of the limonene biosynthetic pathway in these two lineages.
Evolution of plant-microbe interactions
The origin of microbes on Earth, tracing back to the beginning of life more than 3.5 billion years ago, indicates that microbe-microbe interactions have continuously evolved and diversified over time, long before plants started to colonize land 450 million years ago. Therefore, it is likely that both intra- and inter-kingdom intermicrobial interactions represent strong drivers of the establishment of plant-associated microbial consortia at the soil-root interface. Nonetheless, it remains unclear to what extent these interactions in the rhizosphere/phyllosphere and in endophytic plant compartments (i.e., within the host) shape microbial assemblages in nature and whether microbial adaptation to plant habitats drive habitat-specific microbe-microbe interaction strategies that impact plant fitness. Furthermore, the contribution of competitive and cooperative microbe-microbe interactions to the overall community structure remains difficult to evaluate in nature due to the strong environmental noise.
| Biology and health sciences | Basics_4 | Biology |
2238936 | https://en.wikipedia.org/wiki/Projection%20screen | Projection screen | A projection screen is an installation consisting of a surface and a support structure used for displaying a projected image for the view of an audience. Projection screens may be permanently installed on a wall, as in a movie theater, mounted to or placed in a ceiling using a rollable projection surface that retracts into a casing (these can be motorized or manually operated), painted on a wall, or portable with tripod or floor rising models as in a conference room or other non-dedicated viewing space. Another popular type of portable screens are inflatable screens for outdoor movie screening (open-air cinema).
Uniformly white or grey screens are used almost exclusively as to avoid any discoloration to the image, while the most desired brightness of the screen depends on a number of variables, such as the ambient light level and the luminous power of the image source. Flat or curved screens may be used depending on the optics used to project the image and the desired geometrical accuracy of the image production, flat screens being the more common of the two. Screens can be further designed for front or back projection, the more common being front projection systems, which have the image source situated on the same side of the screen as the audience.
Different markets exist for screens targeted for use with digital projectors, movie projectors, overhead projectors and slide projectors, although the basic idea for each of them is very much the same: front projection screens work on diffusely reflecting the light projected on to them, whereas back-projection screens work by diffusely transmitting the light through them.
Screens by installation type in different settings
In the commercial movie theaters, the screen is a reflective surface that may be either aluminized (for high contrast in moderate ambient light) or a white surface with small glass beads (for high brilliance under dark conditions). The screen also has hundreds of small, evenly spaced holes to allow air to and from the speakers and subwoofer, which often are directly behind it.
Rigid wall-mounted screens maintain their geometry perfectly which makes them suitable for applications that demand exact reproduction of image geometry. Such screens are often used in home theaters, along with the pull-down screens.
Pull-down screens (also known as manual wall screens) are often used in spaces where a permanently installed screen would require too much space. These commonly use painted fabric that is rolled in the screen case when not used, making them less obtrusive when the screen is not in use.
Fixed-frame screens provide the greatest level of uniform tension on the screens surface, resulting in the optimal image quality. They are often used in home theater and professional environments where the screen does not need to be recessed into the case.
Electric screens can be wall-mounted, ceiling-mounted or ceiling recessed. These are often larger screens, though electric screens are available for home theater use as well. Electric screens are similar to pull-down screens, but instead of the screen being pulled down manually, an electric motor raises and lowers the screen. Electric screens are usually raised or lowered using either a remote control or wall-mounted switch, although some projectors are equipped with an interface that connects to the screen and automatically lowers the screen when the projector is switched on and raises it when the projector is switched off.
Switchable projection screens can be switched between opaque and clear. In the opaque state, projected image on the screen can be viewed from both sides. It is very good for advertising on store windows.
Mobile screens usually use either a pull-down screen on a free stand, or pull up from a weighted base. These can be used when it is impossible or impractical to mount the screen to a wall or a ceiling.
Both mobile and permanently installed pull-down screens may be of tensioned or not tensioned variety. Tensioned models attempt to keep the fabric flat and immobile, whereas the not tensioned models have the fabric of the screen hanging freely from their support structures. In the latter screens, the fabric can rarely stay immobile if there are currents of air in the room, giving imperfections to the projected image.
Specialty screens may not fall into any of these categories. These include non-solid screens, inflatable screens and others, and can be inexpensively made at home. See the respective articles for more information.
Screen gain
One of the most often quoted properties in a home theater screen is the gain. This is a measure of reflectivity of light compared to a screen coated with magnesium carbonate, titanium dioxide, or
barium sulfate when the measurement is taken for light targeted and reflected perpendicular to the screen. Titanium dioxide is a bright white colour, but greater gains can be accomplished with materials that reflect more of the light parallel to projection axis and less off-axis.
Frequently quoted gain levels of various materials range from 0.8 of light grey matte screens to 2.5 of the more highly reflective glass bead screens. Very high gain levels could be attained simply by using a mirror surface, although the audience would then just see a reflection of the projector, defeating the purpose of using a screen. Many screens with higher gain are simply semi-glossy, and so exhibit more mirror-like properties, namely a bright "hot spot" in the screen—an enlarged (and greatly blurred) reflection of the projector's lens. Opinions differ as to when this "hot spotting" begins to be distracting, but most viewers do not notice differences as large as 30% in the image luminosity, unless presented with a test image and asked to look for variations in brightness. This is possible because humans have greater sensitivity to contrast in smaller details, but less so in luminosity variations as great as half of the screen. Other screens with higher gain are semi-retroreflective. Unlike mirrors, retroreflective surfaces reflect light back toward the source. Hot spotting is less of a problem with retroreflective high-gain screens. At the perpendicular direction used for gain measurement, mirror reflection and retroreflection are indistinguishable, and this has sown confusion about the behavior of high gain screens.
A second common confusion about screen gain arises for grey-colored screens. If a screen material looks grey on casual examination then its total reflectance is much less than 1. However, the grey screen can have measured gain of 1 or even much greater than 1. The geometric behavior of a grey screen is different from that of a white screen of identical gain. Therefore, since geometry is important in screen applications, screen materials should be at least specified by their gain and their total reflectance. Instead of total reflectance, "geometric gain" (equal to the gain divided by the total reflectance) can be the second specification.
Curved screens can be made highly reflective without introducing any visible hot spots, if the curvature of the screen, placement of the projector and the seating arrangement are designed correctly. The object of this design is to have the screen reflect the projected light back to the audience, effectively making the entire screen a giant "hot spot". If the angle of reflection is about the same across the screen, no distracting artifacts will be formed.
Semi-specular high gain screen materials are suited to ceiling-mounted projector setups since the greatest intensity of light will be reflected downward toward the audience at an angle equal and opposite to the angle of incidence. However, for a viewer seated to one side of the audience the opposite side of the screen is much darkened for the same reason. Some structured screen materials are semi-specularly reflective in the vertical plane while more perfectly diffusely reflective in the horizontal plane to avoid this. Glass-bead screens exhibit a phenomenon of retroreflection; the light is reflected more intensely back to its source than in any other direction. They work best for setups where the image source is placed in the same direction from the screen as the audience. With retroreflective screens, the screen center might be brighter than the screen periphery, a kind of hot spotting. This differs from semi-specular screens where the hot spot's location varies depending on the viewer's position in the audience. Retroreflective screens are seen as desirable due to the high image intensity they can produce with a given luminous flux from a projector.
Screen geometry
Projector screens are almost always rectangular in shape. They typically follow a standard display aspect ratio. For most home cinema setups there are two aspect ratios. 16:9 and Cinemascope.
For classroom, businesses and houses of worship settings, 16:10 is the more commonly used projector screen aspect ratio because this matches the aspect ratio used by many modern computers.
Square-shaped screens used for overhead projectors sometimes double as projection screens for digital projectors in meeting rooms, where space is scarce and multiple screens can seem redundant. These screens have an aspect ratio of 1:1 by definition.
Most image sources are designed to project a perfectly rectangular image on a flat screen. If the audience stays relatively close to the projector, a curved screen may be used instead without visible distortion in the image geometry. Viewers closer or farther away will see a pincushion or barrel distortion, and the curved nature of the screen will become apparent when viewed off-axis.
Image brightness and contrast
Apparent contrast in a projected image — the range of brightness — is dependent on the ambient light conditions, luminous power of the projector and the size of the image being projected. A larger screen size means less luminous (luminous power per unit solid angle per unit area) and thus less contrast in the presence of ambient light. Some light will always be created in the room when an image is projected, increasing the ambient light level and thus contributing to the degradation of picture quality. This effect can be lessened by decorating the room with dark colours. The real-room situation is different from the contrast ratios advertised by projector manufacturers, who record the light levels with projector on full black / full white, giving as high contrast ratios as possible.
Manufacturers of home theater screens have attempted to resolve the issue of ambient light by introducing screen surfaces that direct more of the light back to the light source. The rationale behind this approach relies on having the image source placed near the audience, so that the audience will actually see the increased reflected light level on the screen.
Highly reflective flat screens tend to suffer from hot spots, when part of the screen seems much more bright than the rest. This is a result of the high directionality (mirror-likeness) of such screens. Screens with high gain also have a narrower usable viewing angle, as the amount of reflected light rapidly decreases as the viewer moves away from front of such screen. Because of the said effect, these screens are also less vulnerable to ambient light coming from the sides of the screen, as well.
Grey screens
A relatively recent attempt in improving the perceived image quality is the introduction of grey screens, which are more capable of darker tones than their white counterparts. A matte grey screen would have no advantage over a matte white screen in terms of contrast; contemporary grey screens are rather designed to have a gain factor similar to those of matte white screens, but a darker appearance. A darker (grey) screen reflects less light, of course—both light from the projector and ambient light. This decreases the luminance (brightness) of both the projected image and ambient light, so while the light areas of the projected image are dimmer, the dark areas are darker; white is less bright, but intended black is closer to actual black. Many screen manufacturers thus appropriately call their grey screens "high-contrast" models.
Although a projection screen cannot improve a projector's contrast level, the perceived contrast can be boosted.
In an optimal viewing room, the projection screen is reflective, whereas the surroundings are not. The ambient light level is related to the overall reflectivity of the screen, as well as that of the surroundings. In cases where the area of the screen is large compared to that of the surroundings, the screen's contribution to the ambient light may dominate and the effect of the non-screen surfaces of the room may even be negligible. Some examples of this are planetariums and virtual-reality cubes featuring front-projection technology. Some planetariums with dome-shaped projection screens have thus opted to paint the dome interior in gray, in order to reduce the degrading effect of inter-reflections when images of the sun are displayed simultaneously with images of dimmer objects.
Grey screens are designed to rely on powerful image sources that are able to produce adequate levels of luminosity so that the white areas of the image still appear as white, taking advantage of the non-linear perception of brightness in the human eye. People may perceive a wide range of luminosities as "white", as long as the visual clues present in the environment suggest such an interpretation. A grey screen may thus succeed almost as well in delivering a bright-looking image, or fail to do so in other circumstances.
Compared to a white screen, a grey screen reflects less light to the room and less light from the room, making it increasingly effective in dealing with the light originating from the projector. Ambient light originating from other sources may reach the eye immediately after having reflected from the screen surface, giving no advantage over a white high-gain screen in terms of contrast ratio. The potential improvement from a grey screen may thus be best realized in a darkened room, where the only light is that of the projector.
Partly fueled by popularity, grey screen technology has improved greatly in recent years. Grey screens are now available in various gain and grey-scale levels.
Selectively reflective screens
Certain screens are claimed to selectively reflect the narrow wavelengths of projector light while absorbing other wavelengths in the optical spectrum. Sony makes a screen that appears grey in normal room light, and is intended to reduce the effect of ambient light. This is purported to work by preferentially absorbing ambient light of colors not used by the projector, while preferentially reflecting the colors of red, green and blue light the projector uses. A true color-selective screen has not been substantiated. A contrast-enhancing screen has been introduced by Dai Nippon Printing (DNP) and Screen Innovations that is based on thin layers of black louvers rather than wavelength-selective reflection properties.
Screens as an optical element
In an optimally configured system, projection screen surface and the real image plane are made to coincide. From an optical point of view, a screen is not needed for the image to form; screens are rather used to make an image visible.
| Technology | Media and communication: Basics | null |
2240066 | https://en.wikipedia.org/wiki/Vertical%20draft | Vertical draft | In meteorology, an updraft (British English: up-draught) is a small-scale current of rising air, often within a cloud.
Overview
Vertical drafts, known as updrafts or downdrafts, are localized regions of warm or cool air that move vertically. A mass of warm air will typically be less dense than the surrounding region, and so will rise until it reaches air that is either warmer or less dense than itself. The converse will occur for a mass of cool air, and is known as subsidence. This movement of large volumes of air, especially when regions of hot, wet air rise, can create large clouds, and is the central source of thunderstorms. Drafts can also be caused by low or high pressure regions. A low pressure region will attract air from the surrounding area, which will move towards the center and then rise, creating an updraft. A high pressure region will attract air from the surrounding area, which will move towards the center and sink, spawning a downdraft.
Updrafts and downdrafts, along with wind shear in general, are a major contributor to airplane crashes during takeoff and landing in a thunderstorm. Extreme cases, known as downbursts and microbursts, can be deadly and difficult to predict or observe. The crash of Delta Air Lines Flight 191 on its final approach before landing at Dallas/Fort Worth International Airport in 1985 was presumably caused by a microburst, and prompted the Federal Aviation Administration (FAA) to research and deploy new storm detection radar stations at some of the major airports, notably those in the South, Midwest, and Northeast United States where wind shear affects air safety. Downbursts can cause extensive localized damage, similar to that caused by tornadoes. Downburst damage can be differentiated from that of a tornado because the resulting destruction is circular and radiates away from the center. Tornado damage radiates inward, towards the center of the damage.
The term "downdraft" can also refer to a type of backdraft which occurs through chimneys which have fireplaces on the lowermost levels (such as basements) of multi-level buildings. It involves cold air coming down the chimney due to low air pressure, and makes it hard to light fires, and can push soot and carbon monoxide into domiciles.
| Physical sciences | Winds | Earth science |
2241280 | https://en.wikipedia.org/wiki/Veliidae | Veliidae | Veliidae is a family of gregarious predatory insects in the suborder Heteroptera. They are commonly known as riffle bugs, small water striders, or broad-shouldered water striders because the segment immediately behind the head is wider than the rest of the abdomen. Species of the genus Rhagovelia are also referred to as ripple bugs.
Veliidae have a specialized body plan that allows them to walk on water and are neuston. The family Gerridae is another closely related group that is also neuston and both are in the superfamily Gerroidea. Veliidae are smaller however, between . They can be found on ponds, near lake shores, and in rivers worldwide. Some species can also be found on plants near water, in salt water or in mud flats.
Description
Veliidae are very similar to Gerridae. The most consistent characteristic used to separate these two families are internal genitalia differences, however external cues are usually sufficient to tell the families apart.
A general description is as follows: an oval to elongate body covered with hydrofuge hairs. Wings can be present or absent; when present the wings range from well devolved to vestigial. The four segmented antennae is longer than the head and readily visible. The antennae is non-aristate. The eyes are usually large, but there are no ocelli.
Males and females can be differentiated by the fore tibiae. Males have smaller tibiae with a grasping comb, as opposed to the larger plain female tibiae.
Distribution
Veliidae is the largest gerromorphan family and has almost 1173 species and 66 genera. The present distribution of these species points to two centers of origin: one in the Indo-Malayan region and another on the shores of the Caribbean Sea. The geographical distance between these points is probably due to continental drift. And now they are present across all continents (except Antarctica).
Life cycle
Like all Heteroptera, the Veliidae go through an egg, nymph and adult stage. They have four or five nymphal instars. Both the adults and nymphs live together gregariously, in loose communities and can often be found in large groups. Eggs are usually laid underwater, attached to the stream bed, rocks or plant material and held together by a gelatinous substance. In most species females lay under 30 eggs. Nymphs are very similar to adults, but have one segmented tarsus on mid and hind leg as opposed to the adults' two. Some species prefer rapids or riffles in streams but many prefer calmer water.
Behaviour
Veliidae can walk on water because they take advantage of the high surface tension of water and have hydrophobic legs that distribute their weight across more water.
Although Gerridae typically have longer legs, Veliidae also have legs that spread out the weight over a relatively large area. Thousands of hydrofugal hairs also coat the entire body, mitigating potential problems incurred by water contact: air bubbles, trapped among the tiny hairs if the insect is submerged, lift the insect towards the surface again.
Taxonomy
The following genera are recognised in the family Veliidae:
Adriennella
Aegilipsicola
Aegilipsovelia
Aenictovelia
Altavelia
Angilia
Angilovelia
Aphrovelia
Aquulavelia
Arcantivelia
Austromicrovelia
Balticovelia
Baptista
Barbivelia
Brechivelia
Callivelia
Carayonella
Chenevelia
Cylicovelia
Drepanovelia
Electrovelia
Entomovelia
Euvelia
Eyarinella
Fijivelia
Geovelia
Gracilovelia
Haldwania
Halovelia
Haloveloides
Hebrovelia
Husseyella
Lacertovelia
Lathriovelia
Macrovelia
Mangrovelia
Menuthiasia
Microvelia
Microveloidella
Microvelopsis
Millotella
Neoalardus
Nesidovelia
Neusterensifer
Nilsvelia
Ocellovelia
Ocheovelia
Ochthecorisa
Oiovelia
Pacificovelia
Papuavelia
Paravelia
Perittopus
Petrovelia
Phoreticovelia
Platyvelia
Plesiovelia
Polhemovelia
Pseudovelia
Rhagovelia
Rheovelia
Shaverdinia
Starmuhlneria
Steinovelia
Stridulivelia
Strongylovelia
Submicrovelia
Tanyvelia
Tarsovelia
Tarsoveloides
Tenagovelia
Tetraripis
Thirumalaia
Tonkouivelia
Tubuaivelia
Velia
Velohebria
Veloidea
Xenobates
Xiphovelia
Xiphoveloidea
| Biology and health sciences | Hemiptera (true bugs) | Animals |
610813 | https://en.wikipedia.org/wiki/Tarnish | Tarnish | Tarnish is a thin layer of corrosion that forms over copper, brass, aluminum, magnesium, neodymium and other similar metals as their outermost layer undergoes a chemical reaction. Tarnish does not always result from the sole effects of oxygen in the air. For example, silver needs hydrogen sulfide to tarnish, although it may tarnish with oxygen over time. It often appears as a dull, gray or black film or coating over metal. Tarnish is a surface phenomenon that is self-limiting, unlike rust. Only the top few layers of the metal react. The layer of tarnish seals and protects the underlying layers from reacting.
Tarnish preserves the underlying metal in outdoor use, and in this form is called chemical patina. Unlike wear patina necessary in applications such as copper roofing, outdoor copper, bronze, and brass statues and fittings, chemical patina is considered a lot more uneven and undesirable. Patina is the name given to tarnish on copper-based metals, while toning is a term for the type of tarnish which forms on coins.
Chemistry
Tarnish is a product of a chemical reaction between a metal and a nonmetal compound, especially oxygen and sulfur dioxide. It is usually a metal oxide, the product of oxidation; sometimes it is a metal sulfide. The metal oxide sometimes reacts with water to make the hydroxide, or with carbon dioxide to make the carbonate. It is a chemical change. There are various methods to prevent metals from tarnishing.
Prevention and removal
Using a thin coat of polish can prevent tarnish from forming over these metals. Tarnish can be removed by using steel wool, sandpaper, emery paper, baking soda or a file to rub or polish the metal's dull surface. Fine objects (such as silverware) may have the tarnish electrochemically reversed (non-destructively) by resting the objects on a piece of aluminium foil in a pot of boiling water with a small amount of salt or baking soda, or it may be removed with a special polishing compound and a soft cloth. Gentler abrasives, such as calcium carbonate, are often used by museums to clean tarnished silver as they cannot damage or scratch the silver and will not leave unwanted residues.
| Physical sciences | Redox reactions | Chemistry |
611074 | https://en.wikipedia.org/wiki/Point%20mutation | Point mutation | A point mutation is a genetic mutation where a single nucleotide base is changed, inserted or deleted from a DNA or RNA sequence of an organism's genome. Point mutations have a variety of effects on the downstream protein product—consequences that are moderately predictable based upon the specifics of the mutation. These consequences can range from no effect (e.g. synonymous mutations) to deleterious effects (e.g. frameshift mutations), with regard to protein production, composition, and function.
Causes
Point mutations usually take place during DNA replication. DNA replication occurs when one double-stranded DNA molecule creates two single strands of DNA, each of which is a template for the creation of the complementary strand. A single point mutation can change the whole DNA sequence. Changing one purine or pyrimidine may change the amino acid that the nucleotides code for.
Point mutations may arise from spontaneous mutations that occur during DNA replication. The rate of mutation may be increased by mutagens. Mutagens can be physical, such as radiation from UV rays, X-rays or extreme heat, or chemical (molecules that misplace base pairs or disrupt the helical shape of DNA). Mutagens associated with cancers are often studied to learn about cancer and its prevention.
There are multiple ways for point mutations to occur. First, ultraviolet (UV) light and higher-frequency light have ionizing capability, which in turn can affect DNA. Reactive oxygen molecules with free radicals, which are a byproduct of cellular metabolism, can also be very harmful to DNA. These reactants can lead to both single-stranded and double-stranded DNA breaks. Third, bonds in DNA eventually degrade, which creates another problem to keep the integrity of DNA to a high standard. There can also be replication errors that lead to substitution, insertion, or deletion mutations.
Categorization
Transition/transversion categorization
In 1959 Ernst Freese coined the terms "transitions" or "transversions" to categorize different types of point mutations. Transitions are replacement of a purine base with another purine or replacement of a pyrimidine with another pyrimidine. Transversions are replacement of a purine with a pyrimidine or vice versa. There is a systematic difference in mutation rates for transitions (Alpha) and transversions (Beta). Transition mutations are about ten times more common than transversions.
Functional categorization
Nonsense mutations include stop-gain and start-loss. Stop-gain is a mutation that results in a premature termination codon (a stop was gained), which signals the end of translation. This interruption causes the protein to be abnormally shortened. The number of amino acids lost mediates the impact on the protein's functionality and whether it will function whatsoever. Stop-loss is a mutation in the original termination codon (a stop was lost), resulting in abnormal extension of a protein's carboxyl terminus. Start-gain creates an AUG start codon upstream of the original start site. If the new AUG is near the original start site, in-frame within the processed transcript and downstream to a ribosomal binding site, it can be used to initiate translation. The likely effect is additional amino acids added to the amino terminus of the original protein. Frame-shift mutations are also possible in start-gain mutations, but typically do not affect translation of the original protein. Start-loss is a point mutation in a transcript's AUG start codon, resulting in the reduction or elimination of protein production.
Missense mutations code for a different amino acid. A missense mutation changes a codon so that a different protein is created, a non-synonymous change. Conservative mutations result in an amino acid change. However, the properties of the amino acid remain the same (e.g., hydrophobic, hydrophilic, etc.). At times, a change to one amino acid in the protein is not detrimental to the organism as a whole. Most proteins can withstand one or two point mutations before their function changes. Non-conservative mutations result in an amino acid change that has different properties than the wild type. The protein may lose its function, which can result in a disease in the organism. For example, sickle-cell disease is caused by a single point mutation (a missense mutation) in the beta-hemoglobin gene that converts a GAG codon into GUG, which encodes the amino acid valine rather than glutamic acid. The protein may also exhibit a "gain of function" or become activated, such is the case with the mutation changing a valine to glutamic acid in the BRAF gene; this leads to an activation of the RAF protein which causes unlimited proliferative signalling in cancer cells. These are both examples of a non-conservative (missense) mutation.
Silent mutations code for the same amino acid (a "synonymous substitution"). A silent mutation does not affect the functioning of the protein. A single nucleotide can change, but the new codon specifies the same amino acid, resulting in an unmutated protein. This type of change is called synonymous change since the old and new codon code for the same amino acid. This is possible because 64 codons specify only 20 amino acids. Different codons can lead to differential protein expression levels, however.
Single base pair insertions and deletions
Sometimes the term point mutation is used to describe insertions or deletions of a single base pair (which has more of an adverse effect on the synthesized protein due to the nucleotides' still being read in triplets, but in different frames: a mutation called a frameshift mutation).
General consequences
Point mutations that occur in non-coding sequences are most often without consequences, although there are exceptions. If the mutated base pair is in the promoter sequence of a gene, then the expression of the gene may change. Also, if the mutation occurs in the splicing site of an intron, then this may interfere with correct splicing of the transcribed pre-mRNA.
By altering just one amino acid, the entire peptide may change, thereby changing the entire protein. The new protein is called a protein variant. If the original protein functions in cellular reproduction then this single point mutation can change the entire process of cellular reproduction for this organism.
Point germline mutations can lead to beneficial as well as harmful traits or diseases. This leads to adaptations based on the environment where the organism lives. An advantageous mutation can create an advantage for that organism and lead to the trait's being passed down from generation to generation, improving and benefiting the entire population. The scientific theory of evolution is greatly dependent on point mutations in cells. The theory explains the diversity and history of living organisms on Earth. In relation to point mutations, it states that beneficial mutations allow the organism to thrive and reproduce, thereby passing its positively affected mutated genes on to the next generation. On the other hand, harmful mutations cause the organism to die or be less likely to reproduce in a phenomenon known as natural selection.
There are different short-term and long-term effects that can arise from mutations. Smaller ones would be a halting of the cell cycle at numerous points. This means that a codon coding for the amino acid glycine may be changed to a stop codon, causing the proteins that should have been produced to be deformed and unable to complete their intended tasks. Because the mutations can affect the DNA and thus the chromatin, it can prohibit mitosis from occurring due to the lack of a complete chromosome. Problems can also arise during the processes of transcription and replication of DNA. These all prohibit the cell from reproduction and thus lead to the death of the cell. Long-term effects can be a permanent changing of a chromosome, which can lead to a mutation. These mutations can be either beneficial or detrimental. Cancer is an example of how they can be detrimental.
Other effects of point mutations, or single nucleotide polymorphisms in DNA, depend on the location of the mutation within the gene. For example, if the mutation occurs in the region of the gene responsible for coding, the amino acid sequence of the encoded protein may be altered, causing a change in the function, protein localization, stability of the protein or protein complex. Many methods have been proposed to predict the effects of missense mutations on proteins. Machine learning algorithms train their models to distinguish known disease-associated from neutral mutations whereas other methods do not explicitly train their models but almost all methods exploit the evolutionary conservation assuming that changes at conserved positions tend to be more deleterious. While majority of methods provide a binary classification of effects of mutations into damaging and benign, a new level of annotation is needed to offer an explanation of why and how these mutations damage proteins.
Moreover, if the mutation occurs in the region of the gene where transcriptional machinery binds to the protein, the mutation can affect the binding of the transcription factors because the short nucleotide sequences recognized by the transcription factors will be altered. Mutations in this region can affect rate of efficiency of gene transcription, which in turn can alter levels of mRNA and, thus, protein levels in general.
Point mutations can have several effects on the behavior and reproduction of a protein depending on where the mutation occurs in the amino acid sequence of the protein. If the mutation occurs in the region of the gene that is responsible for coding for the protein, the amino acid may be altered. This slight change in the sequence of amino acids can cause a change in the function, activation of the protein meaning how it binds with a given enzyme, where the protein will be located within the cell, or the amount of free energy stored within the protein.
If the mutation occurs in the region of the gene where transcriptional machinery binds to the protein, the mutation can affect the way in which transcription factors bind to the protein. The mechanisms of transcription bind to a protein through recognition of short nucleotide sequences. A mutation in this region may alter these sequences and, thus, change the way the transcription factors bind to the protein. Mutations in this region can affect the efficiency of gene transcription, which controls both the levels of mRNA and overall protein levels.
Specific diseases caused by point mutations
Cancer
Point mutations in multiple tumor suppressor proteins cause cancer. For instance, point mutations in Adenomatous Polyposis Coli promote tumorigenesis. A novel assay, Fast parallel proteolysis (FASTpp), might help swift screening of specific stability defects in individual cancer patients.
Neurofibromatosis
Neurofibromatosis is caused by point mutations in the Neurofibromin 1 or Neurofibromin 2 gene.
Sickle-cell anemia
Sickle-cell anemia is caused by a point mutation in the β-globin chain of hemoglobin, causing the hydrophilic amino acid glutamic acid to be replaced with the hydrophobic amino acid valine at the sixth position.
The β-globin gene is found on the short arm of chromosome 11. The association of two wild-type α-globin subunits with two mutant β-globin subunits forms hemoglobin S (HbS). Under low-oxygen conditions (being at high altitude, for example), the absence of a polar amino acid at position six of the β-globin chain promotes the non-covalent polymerisation (aggregation) of hemoglobin, which distorts red blood cells into a sickle shape and decreases their elasticity.
Hemoglobin is a protein found in red blood cells, and is responsible for the transportation of oxygen through the body. There are two subunits that make up the hemoglobin protein: beta-globins and alpha-globins.
Beta-hemoglobin is created from the genetic information on the HBB, or "hemoglobin, beta" gene found on chromosome 11p15.5. A single point mutation in this polypeptide chain, which is 147 amino acids long, results in the disease known as Sickle Cell Anemia.
Sickle-cell anemia is an autosomal recessive disorder that affects 1 in 500 African Americans, and is one of the most common blood disorders in the United States. The single replacement of the sixth amino acid in the beta-globin, glutamic acid, with valine results in deformed red blood cells. These sickle-shaped cells cannot carry nearly as much oxygen as normal red blood cells and they get caught more easily in the capillaries, cutting off blood supply to vital organs. The single nucleotide change in the beta-globin means that even the smallest of exertions on the part of the carrier results in severe pain and even heart attack. Below is a chart depicting the first thirteen amino acids in the normal and abnormal sickle cell polypeptide chain.
Tay–Sachs disease
The cause of Tay–Sachs disease is a genetic defect that is passed from parent to child. This genetic defect is located in the HEXA gene, which is found on chromosome 15.
The HEXA gene makes part of an enzyme called beta-hexosaminidase A, which plays a critical role in the nervous system. This enzyme helps break down a fatty substance called GM2 ganglioside in nerve cells.
Mutations in the HEXA gene disrupt the activity of beta-hexosaminidase A, preventing the breakdown of the fatty substances. As a result, the fatty substances accumulate to deadly levels in the brain and spinal cord. The buildup of GM2 ganglioside causes progressive damage to the nerve cells. This is the cause of the signs and symptoms of Tay-Sachs disease.
Repeat-induced point mutation
In molecular biology, repeat-induced point mutation or RIP is a process by which DNA accumulates G:C to A:T transition mutations. Genomic evidence indicates that RIP occurs or has occurred in a variety of fungi while experimental evidence indicates that RIP is active in Neurospora crassa, Podospora anserina, Magnaporthe grisea, Leptosphaeria maculans, Gibberella zeae and Nectria haematococca. In Neurospora crassa, sequences mutated by RIP are often methylated de novo.
RIP occurs during the sexual stage in haploid nuclei after fertilization but prior to meiotic DNA replication. In Neurospora crassa, repeat sequences of at least 400 base pairs in length are vulnerable to RIP. Repeats with as low as 80% nucleotide identity may also be subject to RIP. Though the exact mechanism of repeat recognition and mutagenesis are poorly understood, RIP results in repeated sequences undergoing multiple transition mutations.
The RIP mutations do not seem to be limited to repeated sequences. Indeed, for example, in the phytopathogenic fungus L. maculans, RIP mutations are found in single copy regions, adjacent to the repeated elements. These regions are either non-coding regions or genes encoding small secreted proteins including avirulence genes.
The degree of RIP within these single copy regions was proportional to their proximity to repetitive elements.
Rep and Kistler have speculated that the presence of highly repetitive regions containing transposons, may promote mutation of resident effector genes. So the presence of effector genes within such regions is suggested to promote their adaptation and diversification when exposed to strong selection pressure.
As RIP mutation is traditionally observed to be restricted to repetitive regions and not single copy regions, Fudal et al. suggested that leakage of RIP mutation might occur within a relatively short distance of a RIP-affected repeat. Indeed, this has been reported in N. crassa whereby leakage of RIP was detected in single copy sequences at least 930 bp from the boundary of neighbouring duplicated sequences.
To elucidate the mechanism of detection of repeated sequences leading to RIP may allow to understand how the flanking sequences may also be affected.
Mechanism
RIP causes G:C to A:T transition mutations within repeats, however, the mechanism that detects the repeated sequences is unknown. RID is the only known protein essential for RIP. It is a DNA methyltransferease-like protein, that when mutated or knocked out results in loss of RIP. Deletion of the rid homolog in Aspergillus nidulans, dmtA, results in loss of fertility while deletion of the rid homolog in Ascobolus immersens, masc1, results in fertility defects and loss of methylation induced premeiotically (MIP).
Consequences
RIP is believed to have evolved as a defense mechanism against transposable elements, which resemble parasites by invading and multiplying within the genome.
RIP creates multiple missense and nonsense mutations in the coding sequence. This hypermutation of G-C to A-T in repetitive sequences eliminates functional gene products of the sequence (if there were any to begin with). In addition, many of the C-bearing nucleotides become methylated, thus decreasing transcription.
Use in molecular biology
Because RIP is so efficient at detecting and mutating repeats, fungal biologists often use it as a tool for mutagenesis. A second copy of a single-copy gene is first transformed into the genome. The fungus must then mate and go through its sexual cycle to activate the RIP machinery. Many different mutations within the duplicated gene are obtained from even a single fertilization event so that inactivated alleles, usually due to nonsense mutations, as well as alleles containing missense mutations can be obtained.
History
The cellular reproduction process of meiosis was discovered by Oscar Hertwig in 1876. Mitosis was discovered several years later in 1882 by Walther Flemming.
Hertwig studied sea urchins, and noticed that each egg contained one nucleus prior to fertilization and two nuclei after. This discovery proved that one spermatozoon could fertilize an egg, and therefore proved the process of meiosis. Hermann Fol continued Hertwig's research by testing the effects of injecting several spermatozoa into an egg, and found that the process did not work with more than one spermatozoon.
Flemming began his research of cell division starting in 1868. The study of cells was an increasingly popular topic in this time period. By 1873, Schneider had already begun to describe the steps of cell division. Flemming furthered this description in 1874 and 1875 as he explained the steps in more detail. He also argued with Schneider's findings that the nucleus separated into rod-like structures by suggesting that the nucleus actually separated into threads that in turn separated. Flemming concluded that cells replicate through cell division, to be more specific mitosis.
Matthew Meselson and Franklin Stahl are credited with the discovery of DNA replication. Watson and Crick acknowledged that the structure of DNA did indicate that there is some form of replicating process. However, there was not a lot of research done on this aspect of DNA until after Watson and Crick. People considered all possible methods of determining the replication process of DNA, but none were successful until Meselson and Stahl. Meselson and Stahl introduced a heavy isotope into some DNA and traced its distribution. Through this experiment, Meselson and Stahl were able to prove that DNA reproduces semi-conservatively.
| Biology and health sciences | Genetics | Biology |
611177 | https://en.wikipedia.org/wiki/Perchlorate | Perchlorate | A perchlorate is a chemical compound containing the perchlorate ion, , the conjugate base of perchloric acid (ionic perchlorate). As counterions, there can be metal cations, quaternary ammonium cations or other ions, for example, nitronium cation ().
The term perchlorate can also describe perchlorate esters or covalent perchlorates. These are organic compounds that are alkyl or aryl esters of perchloric acid. They are characterized by a covalent bond between an oxygen atom of the ClO4 moiety and an organyl group.
In most ionic perchlorates, the cation is non-coordinating. The majority of ionic perchlorates are commercially produced salts commonly used as oxidizers for pyrotechnic devices and for their ability to control static electricity in food packaging. Additionally, they have been used in rocket propellants, fertilizers, and as bleaching agents in the paper and textile industries.
Perchlorate contamination of food and water endangers human health, primarily affecting the thyroid gland.
Ionic perchlorates are typically colorless solids that exhibit good solubility in water. The perchlorate ion forms when they dissolve in water, dissociating into ions. Many perchlorate salts also exhibit good solubility in non-aqueous solvents. Four perchlorates are of primary commercial interest: ammonium perchlorate , perchloric acid , potassium perchlorate and sodium perchlorate .
Production
Perchlorate salts are typically manufactured through the process of electrolysis, which involves oxidizing aqueous solutions of corresponding chlorates. This technique is commonly employed in the production of sodium perchlorate, which finds widespread use as a key ingredient in rocket fuel. Perchlorate salts are also commonly produced by reacting perchloric acid with bases, such as ammonium hydroxide or sodium hydroxide. Ammonium perchlorate, which is highly valued, can also be produced via an electrochemical process.
Perchlorate esters are formed in the presence of a nucleophilic catalyst via a perchlorate salt's nucleophilic substitution onto an alkylating agent.
Uses
The dominant use of perchlorates is as oxidizers in propellants for rockets, fireworks and highway flares. Of particular value is ammonium perchlorate composite propellant as a component of solid rocket fuel. In a related but smaller application, perchlorates are used extensively within the pyrotechnics industry and in certain munitions and for the manufacture of matches.
Perchlorate is used to control static electricity in food packaging. Sprayed onto containers it stops statically charged food from clinging to plastic or paper/cardboard surface.
Niche uses include lithium perchlorate, which decomposes exothermically to produce oxygen, useful in oxygen "candles" on spacecraft, submarines, and in other situations where a reliable backup oxygen supply is needed.
Potassium perchlorate has, in the past, been used therapeutically to help manage Graves' disease. It impedes production of the thyroid hormones that contain iodine.
As perchlorate is generally a non-complexing anion and that its sodium salts is particularly soluble, it is commonly used as a background, or supporting, electrolyte in solution chemistry, electrophoresis, and electrochemistry. Although used as a powerful oxidizer in propulsive powders and explosives, quite surprisingly, the perchlorate anion is a weak oxidant in aqueous solution because of kinetics limitations severely hindering the electron transfer.
Chemical properties
The perchlorate ion is the least redox reactive of the generalized chlorates. Perchlorate contains chlorine in its highest oxidation number (+7). A table of reduction potentials of the four chlorates shows that, contrary to expectation, perchlorate in aqueous solution is the weakest oxidant among the four.
These data show that the perchlorate and chlorate are stronger oxidizers in acidic conditions than in basic conditions.
Gas phase measurements of heats of reaction (which allow computation of ΔfH°) of various chlorine oxides do follow the expected trend wherein exhibits the largest endothermic value of ΔfH° (238.1 kJ/mol) while exhibits the lowest endothermic value of ΔfH° (80.3 kJ/mol).
Weak base and weak coordinating anion
As perchloric acid is one of the strongest mineral acids, perchlorate is a weak base in the sense of Brønsted–Lowry acid–base theory.
As it is also generally a weakly coordinating anion, perchlorate is commonly used as a background, or supporting, electrolyte.
Weak oxidant in aqueous solution due to kinetic limitations
Perchlorate compounds oxidize organic compounds, especially when the mixture is heated. The explosive decomposition of ammonium perchlorate is catalyzed by metals and heat.
As perchlorate is a weak Lewis base (i.e., a weak electron pair donor) and a weak nucleophilic anion, it is also a very weakly coordinating anion. This is why it is often used as a supporting electrolyte to study the complexation and the chemical speciation of many cations in aqueous solution or in electroanalytical methods (voltammetry, electrophoresis…). Although the perchlorate reduction is thermodynamically favorable , and that is expected to be a strong oxidant, most often in aqueous solution, it is practically an inert species behaving as an extremely slow oxidant because of severe kinetics limitations. The metastable character of perchlorate in the presence of reducing cations such as in solution is due to the difficulty to form an activated complex facilitating the electron transfer and the exchange of oxo groups in the opposite direction. These strongly hydrated cations cannot form a sufficiently stable coordination bridge with one of the four oxo groups of the perchlorate anion. Although thermodynamically a mild reductant, ion exhibits a stronger trend to remain coordinated by water molecules to form the corresponding hexa-aquo complex in solution. The high activation energy of the cation binding with perchlorate to form a transient inner sphere complex more favourable to electron transfer considerably hinders the redox reaction. The redox reaction rate is limited by the formation of a favorable activated complex involving an oxo-bridge between the perchlorate anion and the metallic cation. It depends on the molecular orbital rearrangement (HOMO and LUMO orbitals) necessary for a fast oxygen atom transfer (OAT) and the associated electron transfer as studied experimentally by Henry Taube (1983 Nobel Prize in Chemistry) and theoretically by Rudolph A. Marcus (1992 Nobel Prize in Chemistry), both awarded for their respective works on the mechanisms of electron-transfer reactions with metal complexes and in chemical systems.
In contrast to the cations which remain unoxidized in deaerated perchlorate aqueous solutions free of dissolved oxygen, other cations such as Ru(II) and Ti(III) can form a more stable bridge between the metal centre and one of the oxo groups of . In the inner sphere electron transfer mechanism to observe the perchlorate reduction, the anion must quickly transfer an oxygen atom to the reducing cation. When it is the case, metallic cations can readily reduce perchlorate in solution. Ru(II) can reduce to , while V(II), V(III), Mo(III), Cr(II) and Ti(III) can reduce to .
Some metal complexes, especially those of rhenium, and some metalloenzymes can catalyze the reduction of perchlorate under mild conditions. Perchlorate reductase (see below), a molybdoenzyme, also catalyzes the reduction of perchlorate. Both the Re- and Mo-based catalysts operate via metal-oxo intermediates.
Microbiology
Over 40 phylogenetically and metabolically diverse microorganisms capable of growth using perchlorate as an electron acceptor have been isolated since 1996. Most originate from the Pseudomonadota, but others include the Bacillota, Moorella perchloratireducens and Sporomusa sp., and the archaeon Archaeoglobus fulgidus. With the exception of A. fulgidus, microbes that grow via perchlorate reduction utilize the enzymes perchlorate reductase and chlorite dismutase, which collectively take perchlorate to chloride. In the process, free oxygen () is generated.
Natural abundance
Terrestrial abundance
Perchlorate is created by lightning discharges in the presence of chloride. Perchlorate has been detected in rain and snow samples from Florida and Lubbock, Texas. It is also present in Martian soil.
Naturally occurring perchlorate at its most abundant can be found commingled with deposits of sodium nitrate in the Atacama Desert of northern Chile. These deposits have been heavily mined as sources for nitrate-based fertilizers. Chilean nitrate is in fact estimated to be the source of around of perchlorate imported to the U.S. (1909–1997). Results from surveys of ground water, ice, and relatively unperturbed deserts have been used to estimate a "global inventory" of natural perchlorate presently on Earth.
On Mars
Perchlorate was detected in Martian soil at the level of ~0.6% by weight. It was shown that at the Phoenix landing site it was present as a mixture of 60% and 40% . These salts, formed from perchlorates, act as antifreeze and substantially lower the freezing point of water. Based on the temperature and pressure conditions on present-day Mars at the Phoenix lander site, conditions would allow a perchlorate salt solution to be stable in liquid form for a few hours each day during the summer.
The possibility that the perchlorate was a contaminant brought from Earth was eliminated by several lines of evidence. The Phoenix retro-rockets used ultra pure hydrazine and launch propellants consisting of ammonium perchlorate or ammonium nitrate. Sensors on board Phoenix found no traces of ammonium nitrate, and thus the nitrate in the quantities present in all three soil samples is indigenous to the Martian soil. Perchlorate is widespread in Martian soils at concentrations between 0.5 and 1%. At such concentrations, perchlorate could be an important source of oxygen, but it could also become a critical chemical hazard to astronauts.
In 2006, a mechanism was proposed for the formation of perchlorates that is particularly relevant to the discovery of perchlorate at the Phoenix lander site. It was shown that soils with high concentrations of chloride converted to perchlorate in the presence of titanium dioxide and sunlight/ultraviolet light. The conversion was reproduced in the lab using chloride-rich soils from Death Valley. Other experiments have demonstrated that the formation of perchlorate is associated with wide band gap semiconducting oxides. In 2014, it was shown that perchlorate and chlorate can be produced from chloride minerals under Martian conditions via UV using only NaCl and silicate.
Further findings of perchlorate and chlorate in the Martian meteorite EETA79001 and by the Mars Curiosity rover in 2012-2013 support the notion that perchlorates are globally distributed throughout the Martian surface. With concentrations approaching 0.5% and exceeding toxic levels on Martian soil, Martian perchlorates would present a serious challenge to human settlement, as well as microorganisms. On the other hand, the perchlorate would provide a convenient source of oxygen for the settlements.
On September 28, 2015, NASA announced that analyses of spectral data from the Compact Reconnaissance Imaging Spectrometer for Mars instrument (CRISM) on board the Mars Reconnaissance Orbiter from four different locations where recurring slope lineae (RSL) are present found evidence for hydrated salts. The hydrated salts most consistent with the spectral absorption features are magnesium perchlorate, magnesium chlorate and sodium perchlorate. The findings strongly support the hypothesis that RSL form as a result of contemporary water activity on Mars.
Contamination in environment
Perchlorates are of concern because of uncertainties about toxicity and health effects at low levels in drinking water, impact on ecosystems, and indirect exposure pathways for humans due to accumulation in vegetables. They are water-soluble, exceedingly mobile in aqueous systems, and can persist for many decades under typical groundwater and surface water conditions.
Industrial origin
Perchlorates are used mostly in rocket propellants but also in disinfectants, bleaching agents, and herbicides. Perchlorate contamination is caused during both the manufacture and ignition of rockets and fireworks. Fireworks are also a source of perchlorate in lakes. Removal and recovery methods of these compounds from explosives and rocket propellants include high-pressure water washout, which generates aqueous ammonium perchlorate.
In U.S. drinking water
In 2000, perchlorate contamination beneath the former flare manufacturing plant Olin Corporation Flare Facility, Morgan Hill, California was first discovered several years after the plant had closed. The plant had used potassium perchlorate as one of the ingredients during its 40 years of operation. By late 2003, the State of California and the Santa Clara Valley Water District had confirmed a groundwater plume currently extending over nine miles through residential and agricultural communities.
The California Regional Water Quality Control Board and the Santa Clara Valley Water District have engaged in a major outreach effort, a water well testing program has been underway for about 1,200 residential, municipal, and agricultural wells. Large ion exchange treatment units are operating in three public water supply systems which include seven municipal wells with perchlorate detection. The potentially responsible parties, Olin Corporation and Standard Fuse Incorporated, have been supplying bottled water to nearly 800 households with private wells, and the Regional Water Quality Control Board has been overseeing cleanup efforts.
The source of perchlorate in California was mainly attributed to two manufacturers in the southeast portion of the Las Vegas Valley in Nevada, where perchlorate has been produced for industrial use. This led to perchlorate release into Lake Mead in Nevada and the Colorado River which affected regions of Nevada, California and Arizona, where water from this reservoir is used for consumption, irrigation and recreation for approximately half the population of these states. Lake Mead has been attributed as the source of 90% of the perchlorate in Southern Nevada's drinking water. Based on sampling, perchlorate has been affecting 20 million people, with highest detection in Texas, southern California, New Jersey, and Massachusetts, but intensive sampling of the Great Plains and other middle state regions may lead to revised estimates with additional affected regions. An action level of 18 μg/L has been adopted by several affected states.
In 2001, the chemical was detected at levels as high as 5 μg/L at Joint Base Cape Cod (formerly Massachusetts Military Reservation), over the Massachusetts then state regulation of 2 μg/L.
As of 2009, low levels of perchlorate had been detected in both drinking water and groundwater in 26 states in the U.S., according to the Environmental Protection Agency (EPA).
In food
In 2004, the chemical was found in cow's milk in California at an average level of 1.3 parts per billion (ppb, or μg/L), which may have entered the cows through feeding on crops exposed to water containing perchlorates.
A 2005 study suggested human breast milk had an average of 10.5 μg/L of perchlorate.
From minerals and other natural occurrences
In some places, there is no clear source of perchlorate, and it may be naturally occurring. Natural perchlorate on Earth was first identified in terrestrial nitrate deposits /fertilizers of the Atacama Desert in Chile as early as the 1880s and for a long time considered a unique perchlorate source. The perchlorate released from historic use of Chilean nitrate based fertilizer which the U.S.imported by the hundreds of tons in the early 19th century can still be found in some groundwater sources of the United States, for example Long Island, New York. Recent improvements in analytical sensitivity using ion chromatography based techniques have revealed a more widespread presence of natural perchlorate, particularly in subsoils of Southwest USA, salt evaporites in California and Nevada, Pleistocene groundwater in New Mexico, and even present in extremely remote places such as Antarctica. The data from these studies and others indicate that natural perchlorate is globally deposited on Earth with the subsequent accumulation and transport governed by the local hydrologic conditions.
Despite its importance to environmental contamination, the specific source and processes involved in natural perchlorate production remain poorly understood. Laboratory experiments in conjunction with isotopic studies have implied that perchlorate may be produced on earth by oxidation of chlorine species through pathways involving ozone or its photochemical products. Other studies have suggested that perchlorate can also be formed by lightning activated oxidation of chloride aerosols (e.g., chloride in sea salt sprays), and ultraviolet or thermal oxidation of chlorine (e.g., bleach solutions used in swimming pools) in water.
From nitrate fertilizers
Although perchlorate as an environmental contaminant is usually associated with the manufacture, storage, and testing of solid rocket motors, contamination of perchlorate has been focused as a side effect of the use of natural nitrate fertilizer and its release into ground water. The use of naturally contaminated nitrate fertilizer contributes to the infiltration of perchlorate anions into the ground water and threaten the water supplies of many regions in the US.
One of the main sources of perchlorate contamination from natural nitrate fertilizer use was found to come from the fertilizer derived from Chilean caliche (calcium carbonate), because Chile has rich source of naturally occurring perchlorate anion. Perchlorate concentration was the highest in Chilean nitrate, ranging from 3.3 to 3.98%. Perchlorate in the solid fertilizer ranged from 0.7 to 2.0 mg g−1, variation of less than a factor of 3 and it is estimated that sodium nitrate fertilizers derived from Chilean caliche contain approximately 0.5–2 mg g−1 of perchlorate anion. The direct ecological effect of perchlorate is not well known; its impact can be influenced by factors including rainfall and irrigation, dilution, natural attenuation, soil adsorption, and bioavailability. Quantification of perchlorate concentrations in nitrate fertilizer components via ion chromatography revealed that in horticultural fertilizer components contained perchlorate ranging between 0.1 and 0.46%.
Environmental cleanup
There have been many attempts to eliminate perchlorate contamination. Current remediation technologies for perchlorate have downsides of high costs and difficulty in operation. Thus, there have been interests in developing systems that would offer economic and green alternatives.
Treatment ex situ and in situ
Several technologies can remove perchlorate, via treatments ex situ (away from the location) and in situ (at the location).
Ex situ treatments include ion exchange using perchlorate-selective or nitrite-specific resins, bioremediation using packed-bed or fluidized-bed bioreactors, and membrane technologies via electrodialysis and reverse osmosis. In ex situ treatment via ion exchange, contaminants are attracted and adhere to the ion exchange resin because such resins and ions of contaminants have opposite charge. As the ion of the contaminant adheres to the resin, another charged ion is expelled into the water being treated, in which then ion is exchanged for the contaminant. Ion exchange technology has advantages of being well-suitable for perchlorate treatment and high volume throughput but has a downside that it does not treat chlorinated solvents. In addition, ex situ technology of liquid phase carbon adsorption is employed, where granular activated carbon (GAC) is used to eliminate low levels of perchlorate and pretreatment may be required in arranging GAC for perchlorate elimination.
In situ treatments, such as bioremediation via perchlorate-selective microbes and permeable reactive barrier, are also being used to treat perchlorate. In situ bioremediation has advantages of minimal above-ground infrastructure and its ability to treat chlorinated solvents, perchlorate, nitrate, and RDX simultaneously. However, it has a downside that it may negatively affect secondary water quality. In situ technology of phytoremediation could also be utilized, even though perchlorate phytoremediation mechanism is not fully founded yet.
Bioremediation using perchlorate-reducing bacteria, which reduce perchlorate ions to harmless chloride, has also been proposed.
Health effects
Thyroid inhibition
Perchlorate is a potent competitive inhibitor of the thyroid sodium-iodide symporter. Thus, it has been used to treat hyperthyroidism since the 1950s. At very high doses (70,000–300,000 ppb) the administration of potassium perchlorate was considered the standard of care in the United States, and remains the approved pharmacologic intervention for many countries.
In large amounts perchlorate interferes with iodine uptake into the thyroid gland. In adults, the thyroid gland helps regulate the metabolism by releasing hormones, while in children, the thyroid helps in proper development. The NAS, in its 2005 report, Health Implications of Perchlorate Ingestion, emphasized that this effect, also known as Iodide Uptake Inhibition (IUI) is not an adverse health effect. However, in January 2008, California's Department of Toxic Substances Control stated that perchlorate is becoming a serious threat to human health and water resources. In 2010, the EPA's Office of the Inspector General determined that the agency's own perchlorate reference dose (RfD) of 24.5 parts per billion protects against all human biological effects from exposure, as the federal government is responsible for all US military base groundwater contamination. This finding was due to a significant shift in policy at the EPA in basing its risk assessment on non-adverse effects such as IUI instead of adverse effects. The Office of the Inspector General also found that because the EPA's perchlorate reference dose is conservative and protective of human health further reducing perchlorate exposure below the reference dose does not effectively lower risk.
Because of ammonium perchlorate's adverse effects upon children, Massachusetts set its maximum allowed limit of ammonium perchlorate in drinking water at 2 parts per billion (2 ppb = 2 micrograms per liter).
Perchlorate affects only thyroid hormone. Because it is neither stored nor metabolized, effects of perchlorate on the thyroid gland are reversible, though effects on brain development from lack of thyroid hormone in fetuses, newborns, and children are not.
Toxic effects of perchlorate have been studied in a survey of industrial plant workers who had been exposed to perchlorate, compared to a control group of other industrial plant workers who had no known exposure to perchlorate. After undergoing multiple tests, workers exposed to perchlorate were found to have a significant systolic blood pressure rise compared to the workers who were not exposed to perchlorate, as well as a significant decreased thyroid function compared to the control workers.
A study involving healthy adult volunteers determined that at levels above 0.007 milligrams per kilogram per day (mg/(kg·d)), perchlorate can temporarily inhibit the thyroid gland's ability to absorb iodine from the bloodstream ("iodide uptake inhibition", thus perchlorate is a known goitrogen). The EPA converted this dose into a reference dose of 0.0007 mg/(kg·d) by dividing this level by the standard intraspecies uncertainty factor of 10. The agency then calculated a "drinking water equivalent level" of 24.5 ppb by assuming a person weighs and consumes of drinking water per day over a lifetime.
In 2006, a study reported a statistical association between environmental levels of perchlorate and changes in thyroid hormones of women with low iodine. The study authors were careful to point out that hormone levels in all the study subjects remained within normal ranges. The authors also indicated that they did not originally normalize their findings for creatinine, which would have essentially accounted for fluctuations in the concentrations of one-time urine samples like those used in this study. When the Blount research was re-analyzed with the creatinine adjustment made, the study population limited to women of reproductive age, and results not shown in the original analysis, any remaining association between the results and perchlorate intake disappeared. Soon after the revised Blount Study was released, Robert Utiger, a doctor with the Harvard Institute of Medicine, testified before the US Congress and stated: "I continue to believe that that reference dose, 0.007 milligrams per kilo (24.5 ppb), which includes a factor of 10 to protect those who might be more vulnerable, is quite adequate."
In 2014, a study was published, showing that environmental exposure to perchlorate in pregnant women with hypothyroidism is associated with a significant risk of low IQ in their children.
Lung toxicity
Some studies suggest that perchlorate has pulmonary toxic effects as well. Studies have been performed on rabbits where perchlorate has been injected into the trachea. The lung tissue was removed and analyzed, and it was found that perchlorate injected lung tissue showed several adverse effects when compared to the control group that had been intratracheally injected with saline. Adverse effects included inflammatory infiltrates, alveolar collapse, subpleural thickening, and lymphocyte proliferation.
Aplastic anemia
In the early 1960s, potassium perchlorate used to treat Graves' disease was implicated in the development of aplastic anemia—a condition where the bone marrow fails to produce new blood cells in sufficient quantity—in thirteen patients, seven of whom died. Subsequent investigations have indicated the connection between administration of potassium perchlorate and development of aplastic anemia to be "equivocable at best", which means that the benefit of treatment, if it is the only known treatment, outweighs the risk, and it appeared a contaminant poisoned the 13.
Regulation in the U.S.
Water
In 1998, perchlorate was included in the U.S. EPA Contaminant Candidate List, primarily due to its detection in California drinking water.
In 2002, the EPA completed its draft toxicological review of perchlorate and proposed an reference dose of 0.00003 milligrams per kilogram per day (mg/kg/day) based primarily on studies that identified neurodevelopmental deficits in rat pups. These deficits were linked to maternal exposure to perchlorate.
In 2003, a federal district court in California found that the Comprehensive Environmental Response, Compensation and Liability Act applied, because perchlorate is ignitable, and therefore was a "characteristic" hazardous waste.
Subsequently, the U.S. National Research Council of the National Academy of Sciences (NAS) reviewed the health implications of perchlorate, and in 2005 proposed a much higher reference dose of 0.0007 mg/kg/day based primarily on a 2002 study by Greer et al. During that study, 37 adult human subjects were split into four exposure groups exposed to 0.007 (7 subjects), 0.02 (10 subjects), 0.1 (10 subjects), and 0.5 (10 subjects) mg/kg/day. Significant decreases in iodide uptake were found in the three highest exposure groups. Iodide uptake was not significantly reduced in the lowest exposed group, but four of the seven subjects in this group experienced inhibited iodide uptake. In 2005, the RfD proposed by NAS was accepted by EPA and added to its integrated risk information system (IRIS).
The NAS report described the level of lowest exposure from Greer et al. as a "no-observed-effect level" (NOEL). However, there was actually an effect at that level although not statistically significant largely due to small size of study population (four of seven subjects showed a slight decrease in iodide uptake).
Reduced iodide uptake was not considered to be an adverse effect, even though it is a precursor to an adverse effect, hypothyroidism. Therefore, additional safety factors, would be necessary when extrapolating from the point of departure to the RfD.
Consideration of data uncertainty was insufficient because the Greer, et al. study reflected only a 14-day exposure (=acute) to healthy adults and no additional safety factors were considered to protect sensitive subpopulations like for example, breastfeeding newborns.
Although there has generally been consensus with the Greer et al. study, there has been no consensus with regard to developing a perchlorate RfD. One of the key differences results from how the point of departure is viewed (i.e., NOEL or "lowest-observed-adverse-effect level", LOAEL), or whether a benchmark dose should be used to derive the RfD. Defining the point of departure as a NOEL or LOAEL has implications when it comes to applying appropriate safety factors to the point of departure to derive the RfD.
In early 2006, EPA issued a "Cleanup Guidance" and recommended a Drinking Water Equivalent Level (DWEL) for perchlorate of 24.5 μg/L. Both DWEL and Cleanup Guidance were based on a 2005 review of the existing research by the National Academy of Sciences (NAS).
Lacking a federal drinking water standard, several states subsequently published their own standards for perchlorate including Massachusetts in 2006 and California in 2007. Other states, including Arizona, Maryland, Nevada, New Mexico, New York, and Texas have established non-enforceable, advisory levels for perchlorate.
In 2008, EPA issued an interim drinking water health advisory for perchlorate and with it a guidance and analysis concerning the impacts on the environment and drinking water. California also issued guidance regarding perchlorate use. Both the Department of Defense and some environmental groups voiced questions about the NAS report, but no credible science has emerged to challenge the NAS findings.
In February 2008, the U.S. Food and Drug Administration (FDA) reported that U.S. toddlers on average were being exposed to more than half of EPA's safe dose from food alone. In March 2009, a Centers for Disease Control study found 15 brands of infant formula contaminated with perchlorate and that combined with existing perchlorate drinking water contamination, infants could be at risk for perchlorate exposure above the levels considered safe by EPA.
In 2010, the Massachusetts Department of Environmental Protection set a 10 fold lower RfD (0.07 μg/kg/day) than the NAS RfD using a much higher uncertainty factor of 100. They also calculated an Infant drinking water value, which neither US EPA nor CalEPA had done.
On February 11, 2011, EPA determined that perchlorate meets the Safe Drinking Water Act criteria for regulation as a contaminant. The agency found that perchlorate may have an adverse effect on the health of persons and is known to occur in public water systems with a frequency and at levels that it presents a public health concern. Since then EPA has continued to determine what level of contamination is appropriate. EPA prepared extensive responses to submitted public comments.
In 2016, the Natural Resources Defense Council (NRDC) filed a lawsuit to accelerate EPA's regulation of perchlorate.
In 2019, EPA proposed a Maximum Contaminant Level of 0.056 mg/L for public water systems.
On June 18, 2020, EPA announced that it was withdrawing its 2011 regulatory determination and its 2019 proposal, stating that it had taken "proactive steps" with state and local governments to address perchlorate contamination. In September 2020 NRDC filed suit against EPA for its failure to regulate perchlorate, and stated that 26 million people may be affected by perchlorate in their drinking water. On March 31, 2022, the EPA announced that a review confirmed its 2020 decision. Following the NRDC lawsuit, in 2023 the US Court of Appeals for the DC Circuit ordered EPA to develop a perchlorate standard for public water systems. EPA stated that it will publish a proposed standard for perchlorate in 2025, and issue a final rule in 2027.
Covalent perchlorates
Although typically found as a non-coordinating anion, a few metal complexes are known. Hexaperchloratoaluminate and tetraperchloratoaluminate are strong oxidising agents.
Several perchlorate esters are known. For example, methyl perchlorate is a high energy material that is a strong alkylating agent. Chlorine perchlorate is a covalent inorganic analog.
Safety
As discussed above, iodide is competitor in the thyroid glads. In the presence of reductants, perchlorate forms potentially explosive mixtures. The PEPCON disaster destroyed a production plant for ammonium perchlorate when a fire caused the ammonium perchlorate stored on site to react with the aluminum that the storage tanks were constructed with and explode.
| Physical sciences | Halide oxyanions | Chemistry |
611229 | https://en.wikipedia.org/wiki/Dissection | Dissection | Dissection (from Latin "to cut to pieces"; also called anatomization) is the dismembering of the body of a deceased animal or plant to study its anatomical structure. Autopsy is used in pathology and forensic medicine to determine the cause of death in humans. Less extensive dissection of plants and smaller animals preserved in a formaldehyde solution is typically carried out or demonstrated in biology and natural science classes in middle school and high school, while extensive dissections of cadavers of adults and children, both fresh and preserved are carried out by medical students in medical schools as a part of the teaching in subjects such as anatomy, pathology and forensic medicine. Consequently, dissection is typically conducted in a morgue or in an anatomy lab.
Dissection has been used for centuries to explore anatomy. Objections to the use of cadavers have led to the use of alternatives including virtual dissection of computer models.
In the field of surgery, the term "dissection" or "dissecting" means more specifically the practice of separating an anatomical structure (an organ, nerve or blood vessel) from its surrounding connective tissue in order to minimize unwanted damage during a surgical procedure.
Overview
Plant and animal bodies are dissected to analyze the structure and function of its components. Dissection is practised by students in courses of biology, botany, zoology, and veterinary science, and sometimes in arts studies. In medical schools, students dissect human cadavers to learn anatomy. Zoötomy is sometimes used to describe "dissection of an animal".
Human dissection
A key principle in the dissection of human cadavers (sometimes called androtomy) is the prevention of human disease to the dissector. Prevention of transmission includes the wearing of protective gear, ensuring the environment is clean, dissection technique and pre-dissection tests to specimens for the presence of HIV and hepatitis viruses. Specimens are dissected in morgues or anatomy labs. When provided, they are evaluated for use as a "fresh" or "prepared" specimen. A "fresh" specimen may be dissected within some days, retaining the characteristics of a living specimen, for the purposes of training. A "prepared" specimen may be preserved in solutions such as formalin and pre-dissected by an experienced anatomist, sometimes with the help of a diener. This preparation is sometimes called prosection.
Most dissection involves the careful isolation and removal of individual organs, called the Virchow technique. An alternative more cumbersome technique involves the removal of the entire organ body, called the Letulle technique. This technique allows a body to be sent to a funeral director without waiting for the sometimes time-consuming dissection of individual organs. The Rokitansky method involves an in situ dissection of the organ block, and the technique of Ghon involves dissection of three separate blocks of organs - the thorax and cervical areas, gastrointestinal and abdominal organs, and urogenital organs. Dissection of individual organs involves accessing the area in which the organ is situated, and systematically removing the anatomical connections of that organ to its surroundings. For example, when removing the heart, connects such as the superior vena cava and inferior vena cava are separated. If pathological connections exist, such as a fibrous pericardium, then this may be deliberately dissected along with the organ.
Autopsy and necropsy
Dissection is used to help to determine the cause of death in autopsy (called necropsy in other animals) and is an intrinsic part of forensic medicine.
History
Classical antiquity
Human dissections were carried out by the Greek physicians Herophilus of Chalcedon and Erasistratus of Chios in the early part of the third century BC. Before then, animal dissection had been carried out systematically starting from the fifth century BC. During this period, the first exploration into full human anatomy was performed rather than a base knowledge gained from 'problem-solution' delving. While there was a deep taboo in Greek culture concerning human dissection, there was at the time a strong push by the Ptolemaic government to build Alexandria into a hub of scientific study. For a time, Roman law forbade dissection and autopsy of the human body, so anatomists relied on the cadavers of animals or made observations of human anatomy from injuries of the living. Galen, for example, dissected the Barbary macaque and other primates, assuming their anatomy was basically the same as that of humans, and supplemented these observations with knowledge of human anatomy which he acquired while tending to wounded gladiators.
Celsus wrote in On Medicine I Proem 23, "Herophilus and Erasistratus proceeded in by far the best way: they cut open living men - criminals they obtained out of prison from the kings and they observed, while their subjects still breathed, parts that nature had previously hidden, their position, color, shape, size, arrangement, hardness, softness, smoothness, points of contact, and finally the processes and recesses of each and whether any part is inserted into another or receives the part of another into itself."
Galen was another such writer who was familiar with the studies of Herophilus and Erasistratus.
India
The ancient societies that were rooted in India left behind artwork on how to kill animals during a hunt. The images showing how to kill most effectively depending on the game being hunted relay an intimate knowledge of both external and internal anatomy as well as the relative importance of organs. The knowledge was mostly gained through hunters preparing the recently captured prey. Once the roaming lifestyle was no longer necessary it was replaced in part by the civilization that formed in the Indus Valley. Unfortunately, there is little that remains from this time to indicate whether or not dissection occurred, the civilization was lost to the Aryan people migrating.
Early in the history of India (2nd to 3rd century), the Arthashastra described the 4 ways that death can occur and their symptoms: drowning, hanging, strangling, or asphyxiation. According to that source, an autopsy should be performed in any case of untimely demise.
The practice of dissection flourished during the 7th and 8th century. It was under their rule that medical education was standardized. This created a need to better understand human anatomy, so as to have educated surgeons. Dissection was limited by the religious taboo on cutting the human body. This changed the approach taken to accomplish the goal. The process involved the loosening of the tissues in streams of water before the outer layers were sloughed off with soft implements to reach the musculature. To perfect the technique of slicing, the prospective students used gourds and squash. These techniques of dissection gave rise to an advanced understanding of the anatomy and the enabled them to complete procedures used today, such as rhinoplasty.
During medieval times the anatomical teachings from India spread throughout the known world; however, the practice of dissection was stunted by Islam. The practice of dissection at a university level was not seen again until 1827, when it was performed by the student Pandit Madhusudan Gupta. Through the 1900s, the university teachers had to continually push against the social taboos of dissection, until around 1850 when the universities decided that it was more cost effective to train Indian doctors than bring them in from Britain. Indian medical schools were, however, training female doctors well before those in England.
The current state of dissection in India is deteriorating. The number of hours spent in dissection labs during medical school has decreased substantially over the last twenty years. The future of anatomy education will probably be an elegant mix of traditional methods and integrative computer learning. The use of dissection in early stages of medical training has been shown more effective in the retention of the intended information than their simulated counterparts. However, there is use for the computer-generated experience as review in the later stages. The combination of these methods is intended to strengthen the students' understanding and confidence of anatomy, a subject that is infamously difficult to master. There is a growing need for anatomist—seeing as most anatomy labs are taught by graduates hoping to complete degrees in anatomy—to continue the long tradition of anatomy education.
Islamic world
From the beginning of the Islamic faith in 610 A.D., Shari'ah law has applied to a greater or lesser extent within Muslim countries, supported by Islamic scholars such as Al-Ghazali. Islamic physicians such as Ibn Zuhr (Avenzoar) (1091–1161) in Al-Andalus, Saladin's physician Ibn Jumay during the 12th century, Abd el-Latif in Egypt , and Ibn al-Nafis in Syria and Egypt in the 13th century may have practiced dissection, but it remains ambiguous whether or not human dissection was practiced. Ibn al-Nafis, a physician and Muslim jurist, suggested that the "precepts of Islamic law have discouraged us from the practice of dissection, along with whatever compassion is in our temperament", indicating that while there was no law against it, it was nevertheless uncommon. Islam dictates that the body be buried as soon as possible, barring religious holidays, and that there be no other means of disposal such as cremation. Prior to the 10th century, dissection was not performed on human cadavers. The book Al-Tasrif, written by Al-Zahrawi in 1000 A.D., details surgical procedure that differed from the previous standards. The book was an educational text of medicine and surgery which included detailed illustrations. It was later translated and took the place of Avicenna's The Canon of Medicine as the primary teaching tool in Europe from the 12th century to the 17th century. There were some that were willing to dissect humans up to the 12th century, for the sake of learning, after which it was forbidden. This attitude remained constant until 1952, when the Islamic School of Jurisprudence in Egypt ruled that "necessity permits the forbidden". This decision allowed for the investigation of questionable deaths by autopsy. In 1982, the decision was made by a fatwa that if it serves justice, autopsy is worth the disadvantages. Though Islam now approves of autopsy, the Islamic public still disapproves. Autopsy is prevalent in most Muslim countries for medical and judicial purposes. In Egypt it holds an important place within the judicial structure, and is taught at all the country's medical universities. In Saudi Arabia, whose law is completely dictated by Shari'ah, autopsy is viewed poorly by the population but can be compelled in criminal cases; human dissection is sometimes found at university level. Autopsy is performed for judicial purposes in Qatar and Tunisia. Human dissection is present in the modern day Islamic world, but is rarely published on due to the religious and social stigma.
Tibet
Tibetan medicine developed a rather sophisticated knowledge of anatomy, acquired from long-standing experience with human dissection. Tibetans had adopted the practice of sky burial because of the country's hard ground, frozen for most of the year, and the lack of wood for cremation. A sky burial begins with a ritual dissection of the deceased, and is followed by the feeding of the parts to vultures on the hill tops. Over time, Tibetan anatomical knowledge found its way into Ayurveda and to a lesser extent into Chinese medicine.
Christian Europe
Throughout the history of Christian Europe, the dissection of human cadavers for medical education has experienced various cycles of legalization and proscription in different countries. Dissection was rare during the Middle Ages, but it was practised, with evidence from at least as early as the 13th century. The practice of autopsy in Medieval Western Europe is "very poorly known" as few surgical texts or conserved human dissections have survived.
A modern Jesuit scholar has claimed that the Christian theology contributed significantly to the revival of human dissection and autopsy by providing a new socio-religious and cultural context in which the human cadaver was no longer seen as sacrosanct.
A non-existent edict of the 1163 Council of Tours and an early 14th-century decree of Pope Boniface VIII have mistakenly been identified as prohibiting dissection and autopsy; misunderstanding or extrapolation from these edicts may have contributed to reluctance to perform such procedures. The Middle Ages witnessed the revival of an interest in medical studies, including human dissection and autopsy.
Frederick II (1194–1250), the Holy Roman Emperor, decreed that any that were studying to be a physician or a surgeon must attend a human dissection, which would be held no less than every five years. Some European countries began legalizing the dissection of executed criminals for educational purposes in the late 13th and early 14th centuries. Mondino de Luzzi carried out the first recorded public dissection around 1315. At this time, autopsies were carried out by a team consisting of a Lector, who lectured; the Sector, who did the dissection; and the Ostensor, who pointed to features of interest.
The Italian Galeazzo di Santa Sofia made the first public dissection north of the Alps in Vienna in 1404.
Vesalius in the 16th century carried out numerous dissections in his extensive anatomical investigations. He was attacked frequently for his disagreement with Galen's opinions on human anatomy. Vesalius was the first to lecture and dissect the cadaver simultaneously.
The Catholic Church is known to have ordered an autopsy on conjoined twins Joana and Melchiora Ballestero in Hispaniola in 1533 to determine whether they shared a soul. They found that there were two distinct hearts, and hence two souls, based on the ancient Greek philosopher Empedocles, who believed the soul resided in the heart.
Human dissection was also practised by Renaissance artists. Though most chose to focus on the external surfaces of the body, some like Michelangelo Buonarotti, Antonio del Pollaiuolo, Baccio Bandinelli, and Leonardo da Vinci sought a deeper understanding. However, there were no provisions for artists to obtain cadavers, so they had to resort to unauthorised means, as indeed anatomists sometimes did, such as grave robbing, body snatching, and murder.
Anatomization was sometimes ordered as a form of punishment, as, for example, in 1806 to James Halligan and Dominic Daley after their public hanging in Northampton, Massachusetts.
In modern Europe, dissection is routinely practised in biological research and education, in medical schools, and to determine the cause of death in autopsy. It is generally considered a necessary part of learning and is thus accepted culturally. It sometimes attracts controversy, as when Odense Zoo decided to dissect lion cadavers in public before a "self-selected audience".
Britain
In Britain, dissection remained entirely prohibited from the end of the Roman conquest and through the Middle Ages to the 16th century, when a series of royal edicts gave specific groups of physicians and surgeons some limited rights to dissect cadavers. The permission was quite limited: by the mid-18th century, the Royal College of Physicians and Company of Barber-Surgeons were the only two groups permitted to carry out dissections, and had an annual quota of ten cadavers between them. As a result of pressure from anatomists, especially in the rapidly growing medical schools, the Murder Act 1752 allowed the bodies of executed murderers to be dissected for anatomical research and education. By the 19th century this supply of cadavers proved insufficient, as the public medical schools were growing, and the private medical schools lacked legal access to cadavers. A thriving black market arose in cadavers and body parts, leading to the creation of the profession of body snatching, and the infamous Burke and Hare murders in 1828, when 16 people were murdered for their cadavers, to be sold to anatomists. The resulting public outcry led to the passage of the Anatomy Act 1832, which increased the legal supply of cadavers for dissection.
By the 21st century, the availability of interactive computer programs and changing public sentiment led to renewed debate on the use of cadavers in medical education. The Peninsula College of Medicine and Dentistry in the UK, founded in 2000, became the first modern medical school to carry out its anatomy education without dissection.
United States
In the United States, dissection of frogs became common in college biology classes from the 1920s, and were gradually introduced at earlier stages of education. By 1988, some 75 to 80 percent of American high school biology students were participating in a frog dissection, with a trend towards introduction in elementary schools. The frogs are most commonly from the genus Rana. Other popular animals for high-school dissection at the time of that survey were, among vertebrates, fetal pigs, perch, and cats; and among invertebrates, earthworms, grasshoppers, crayfish, and starfish. About six million animals are dissected each year in United States high schools (2016), not counting medical training and research. Most of these are purchased already dead from slaughterhouses and farms.
Dissection in U.S. high schools became prominent in 1987, when a California student, Jenifer Graham, sued to require her school to let her complete an alternative project. The court ruled that mandatory dissections were permissible, but that Graham could ask to dissect a frog that had died of natural causes rather than one that was killed for the purposes of dissection; the practical impossibility of procuring a frog that had died of natural causes in effect let Graham opt out of the required dissection. The suit gave publicity to anti-dissection advocates. Graham appeared in a 1987 Apple Computer commercial for the virtual-dissection software Operation Frog. The state of California passed a Student's Rights Bill in 1988 requiring that objecting students be allowed to complete alternative projects. Opting out of dissection increased through the 1990s.
In the United States, 17 states along with Washington, D.C. have enacted dissection-choice laws or policies that allow students in primary and secondary education to opt out of dissection. Other states including Arizona, Hawaii, Minnesota, Texas, and Utah have more general policies on opting out on moral, religious, or ethical grounds. To overcome these concerns, J. W. Mitchell High School in New Port Richey, Florida, in 2019 became the first US high school to use synthetic frogs for dissection in its science classes, instead of preserved real frogs.
As for the dissection of cadavers in undergraduate and medical school, traditional dissection is supported by professors and students, with some opposition, limiting the availability of dissection. Upper-level students who have experienced this method along with their professors agree that "Studying human anatomy with colorful charts is one thing. Using a scalpel and an actual, recently-living person is an entirely different matter."
Acquisition of cadavers
The way in which cadaveric specimens are obtained differs greatly according to country. In the UK, donation of a cadaver is wholly voluntary. Involuntary donation plays a role in about 20 percent of specimens in the US and almost all specimens donated in some countries such as South Africa and Zimbabwe. Countries that practice involuntary donation may make available the bodies of dead criminals or unclaimed or unidentified bodies for the purposes of dissection. Such practices may lead to a greater proportion of the poor, homeless and social outcasts being involuntarily donated. Cadavers donated in one jurisdiction may also be used for the purposes of dissection in another, whether across states in the US, or imported from other countries, such as with Libya. As an example of how a cadaver is donated voluntarily, a funeral home in conjunction with a voluntary donation program identifies a body who is part of the program. After broaching the subject with relatives in a diplomatic fashion, the body is then transported to a registered facility. The body is tested for the presence of HIV and hepatitis viruses. It is then evaluated for use as a "fresh" or "prepared" specimen.
Disposal of specimens
Cadaveric specimens for dissection are, in general, disposed of by cremation. The deceased may then be interred at a local cemetery. If the family wishes, the ashes of the deceased are then returned to the family. Many institutes have local policies to engage, support and celebrate the donors. This may include the setting up of local monuments at the cemetery.
Use in education
Human cadavers are often used in medicine to teach anatomy or surgical instruction. Cadavers are selected according to their anatomy and availability. They may be used as part of dissection courses involving a "fresh" specimen so as to be as realistic as possible—for example, when training surgeons. Cadavers may also be pre-dissected by trained instructors. This form of dissection involves the preparation and preservation of specimens for a longer time period and is generally used for the teaching of anatomy.
Alternatives
Some alternatives to dissection may present educational advantages over the use of animal cadavers, while eliminating perceived ethical issues. These alternatives include computer programs, lectures, three dimensional models, films, and other forms of technology. Concern for animal welfare is often at the root of objections to animal dissection. Studies show that some students reluctantly participate in animal dissection out of fear of real or perceived punishment or ostracism from their teachers and peers, and many do not speak up about their ethical objections.
One alternative to the use of cadavers is computer technology. At Stanford Medical School, software combines X-ray, ultrasound and MRI imaging for display on a screen as large as a body on a table. In a variant of this, a "virtual anatomy" approach being developed at New York University, students wear three dimensional glasses and can use a pointing device to "[swoop] through the virtual body, its sections as brightly colored as living tissue." This method is claimed to be "as dynamic as Imax [cinema]".
Advantages and disadvantages
Proponents of animal-free teaching methodologies argue that alternatives to animal dissection can benefit educators by increasing teaching efficiency and lowering instruction costs while affording teachers an enhanced potential for the customization and repeat-ability of teaching exercises. Those in favor of dissection alternatives point to studies which have shown that computer-based teaching methods "saved academic and nonacademic staff time ... were considered to be less expensive and an effective and enjoyable mode of student learning [and] ... contributed to a significant reduction in animal use" because there is no set-up or clean-up time, no obligatory safety lessons, and no monitoring of misbehavior with animal cadavers, scissors, and scalpels.
With software and other non-animal methods, there is also no expensive disposal of equipment or hazardous material removal. Some programs also allow educators to customize lessons and include built-in test and quiz modules that can track student performance. Furthermore, animals (whether dead or alive) can be used only once, while non-animal resources can be used for many years—an added benefit that could result in significant cost savings for teachers, school districts, and state educational systems.
Several peer-reviewed comparative studies examining information retention and performance of students who dissected animals and those who used an alternative instruction method have concluded that the educational outcomes of students who are taught basic and advanced biomedical concepts and skills using non-animal methods are equivalent or superior to those of their peers who use animal-based laboratories such as animal dissection.
Some reports state that students' confidence, satisfaction, and ability to retrieve and communicate information was much higher for those who participated in alternative activities compared to dissection. Three separate studies at universities across the United States found that students who modeled body systems out of clay were significantly better at identifying the constituent parts of human anatomy than their classmates who performed animal dissection.
Another study found that students preferred using clay modeling over animal dissection and performed just as well as their cohorts who dissected animals.
In 2008, the National Association of Biology Teachers (NABT) affirmed its support for classroom animal dissection stating that they "Encourage the presence of live animals in the classroom with appropriate consideration to the age and maturity level of the students ... NABT urges teachers to be aware that alternatives to dissection have their limitations. NABT supports the use of these materials as adjuncts to the educational process but not as exclusive replacements for the use of actual organisms."
The National Science Teachers Association (NSTA) "supports including live animals as part of instruction in the K-12 science classroom because observing and working with animals firsthand can spark students' interest in science as well as a general respect for life while reinforcing key concepts" of biological sciences. NSTA also supports offering dissection alternatives to students who object to the practice.
The NORINA database lists over 3,000 products which may be used as alternatives or supplements to animal use in education and training. These include alternatives to dissection in schools. InterNICHE has a similar database and a loans system.
Additional images
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611249 | https://en.wikipedia.org/wiki/Bell%20tower | Bell tower | A bell tower is a tower that contains one or more bells, or that is designed to hold bells even if it has none. Such a tower commonly serves as part of a Christian church, and will contain church bells, but there are also many secular bell towers, often part of a municipal building, an educational establishment, or a tower built specifically to house a carillon. Church bell towers often incorporate clocks, and secular towers usually do, as a public service.
The term campanile (, also , ), from the Italian campanile, which in turn derives from campana, meaning "bell", is synonymous with bell tower; though in English usage campanile tends to be used to refer to a free standing bell tower. A bell tower may also in some traditions be called a belfry, though this term may also refer specifically to the substructure that houses the bells and the ringers rather than the complete tower.
The tallest free-standing bell tower in the world, high, is the Mortegliano Bell Tower, in the Friuli-Venezia Giulia region, Italy.
Purpose
Bells are rung from a tower to enable them to be heard at a distance. Church bells can signify the time for worshippers to go to church for a communal service, and can be an indication of the fixed times of daily Christian prayer, called the canonical hours, which number seven and are contained in breviaries. They are also rung on special occasions such as a wedding, or a funeral service. In some religious traditions they are used within the liturgy of the church service to signify to people that a particular part of the service has been reached.
A bell tower may have a single bell, or a collection of bells which are tuned to a common scale. They may be stationary and chimed, rung randomly by swinging through a small arc, or swung through a full circle to enable the high degree of control of English change ringing. They may house a carillon or chimes, in which the bells are sounded by hammers connected via cables to a keyboard. These can be found in many churches and secular buildings in Europe and America including college and university campuses.
A variety of electronic devices exist to simulate the sound of bells, but any substantial tower in which a considerable sum of money has been invested will generally have a real set of bells.
Some churches have an exconjuratory in the bell tower, a space where ceremonies were conducted to ward off weather-related calamities, like storms and excessive rain. The main bell tower of the Cathedral of Murcia has four.
In Christianity, many churches ring their church bells from belltowers three times a day, at 9 am, 12 pm and 3 pm to summon the Christian faithful to recite the Lord's Prayer; the injunction to pray the Lord's prayer thrice daily was given in Didache 8, 2 f., which, in turn, was influenced by the Jewish practice of praying thrice daily found in the Old Testament, specifically in , which suggests "evening and morning and at noon", and , in which the prophet Daniel prays thrice a day. The early Christians thus came to pray the Lord's Prayer at 9 am, 12 pm and 3 pm; as such, in Christianity, many Lutheran and Anglican churches ring their church bells from belltowers three times a day: in the morning, at noon and in the evening calling Christians to recite the Lord's Prayer. Many Catholic Christian churches ring their bells thrice a day, at 6a.m., noon, and 6p.m., to call the faithful to recite the Angelus, a prayer recited in honour of the Incarnation of God. Oriental Orthodox Christians, such as Copts and Indians, use a breviary such as the Agpeya and Shehimo to pray the canonical hours seven times a day while facing in the eastward direction; church bells are tolled, especially in monasteries, to mark these seven fixed prayer times (cf. ).
The Christian tradition of the ringing of church bells from a belltower is analogous to Islamic tradition of the adhan (call to prayer) from a minaret.
Old bell towers which are no longer used for their original purpose may be kept for their historic or architectural value, though in countries with a strong campanological tradition they often continue to have the bells rung.
History
Europe
In 400 AD, Paulinus of Nola introduced church bells into the Christian Church. By the 11th century, bells housed in belltowers became commonplace.
Historic bell towers exist throughout Europe. The Irish round towers are thought to have functioned in part as bell towers. Famous medieval European examples include Bruges (Belfry of Bruges), Ypres (Cloth Hall, Ypres), Ghent (Belfry of Ghent). Perhaps the most famous European free-standing bell tower, however, is the so-called "Leaning Tower of Pisa", which is the campanile of the Duomo di Pisa in Pisa, Italy. In 1999 thirty-two Belgian belfries were added to the UNESCO's list of World Heritage Sites. In 2005 this list was extended with one Belgian and twenty-three Northern French belfries and is since known as Belfries of Belgium and France. Most of these were attached to civil buildings, mainly city halls, as symbols of the greater power the cities in the region got in the Middle Ages; a small number of buildings not connected with a belfry, such as bell towers of—or with their—churches, also occur on this same list (details). In the Middle Ages, cities sometimes kept their important documents in belfries. Not all are on a large scale; the "bell" tower of Katúň, in Slovakia, is typical of the many more modest structures that were once common in country areas. Archaic wooden bell towers survive adjoining churches in Lithuania and as well as in some parts of Poland.
In Orthodox Eastern Europe bell ringing also has a strong cultural significance (Russian Orthodox bell ringing), and churches were constructed with bell towers (see also List of tall Orthodox Bell towers).
China
Bell towers (Chinese: Zhonglou, Japanese: Shōrō) are common in China and the countries of related cultures. They may appear both as part of a temple complex and as an independent civic building, often paired with a drum tower, as well as in local church buildings. Among the best known examples are the Bell Tower (Zhonglou) of Beijing and the Bell Tower of Xi'an.
Gallery
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611501 | https://en.wikipedia.org/wiki/Coral%20snake | Coral snake | Coral snakes are a large group of elapid snakes that can be divided into two distinct groups, the Old World coral snakes and New World coral snakes. There are 27 species of Old World coral snakes, in three genera (Calliophis, Hemibungarus, and Sinomicrurus), and 83 recognized species of New World coral snakes, in two genera (Micruroides and Micrurus). Genetic studies have found that the most basal lineages have origins in Asia, suggesting that the group originated in the Old World. While new world species of both genera are venomous, their bites are seldom lethal; only two confirmed fatalities have been documented in the past 100 years from the genus Micrurus. Meanwhile, snakes of the genus Micruroides have never caused a medically significant bite.
North American coloration patterns
Experts now recognize that certain coloration patterns and common mnemonics—such as the phrase “Red against yellow, deadly fellow; red against black, friendly Jack,” which people sometimes use to distinguish between the venomous coral snake and the non-venomous milksnake—are not consistent enough to be trustworthy. While any snake exhibiting the coral snake's color and/or banding pattern in the southeastern United States will almost certainly, in fact, be a coral snake, there are coral snakes in other parts of the world that are colored differently.
Coral snakes in the United States are most notable for their red, yellow/white, and black-colored banding. However, several nonvenomous species in the U.S. have similar (though not identical) bandings, including the two scarlet snake species in the genus Cemophora, and some of the kingsnakes (including the aforementioned milksnakes) in the genus Lampropeltis. However, in reference to the mnemonic phrase “red touching yellow, a deadly fellow,” some of these kingsnakes do not naturally display any red touching yellow, to begin with. Additionally, some ground snakes in the genus Sonora (of the southwestern U.S.) can have a color pattern that matches that of the sympatric Sonoran coral snake (Micruroides euryxanthus). No genuine coral snake in the U.S. exhibits red bands of color, in contact with bands of black, except in rare cases of an aberrant pattern. Thus, while on extremely rare occasions when a certain non-venomous snake might be mistaken for a coral snake, the mnemonic holds true. However, a red–yellow–black banded snake in the U.S. (whose red and black banding actually touch) is rarely a venomous coral snake.
Furthermore, the mnemonic is not consistently accurate for North American coral snake species found south of the U.S., either. Some species, like Mexico’s Oaxacan coral snake (Micrurus ephippifer) or Clark's coral snake (Micrurus clarki, of Costa Rica & Panama), do actually fit the mnemonic; others, like the Honduran Roatan coral snake (Micrurus ruatanus), or the redtail coral snake (Micrurus mipartitus) of Panama, do not. Further still, some South American coral snakes will fit the mnemonic, while others do not. To complicate the issue more, the South American tricolored hognose snake (Xenodon pulcher) has repeated bands of red-black-white-black, imitating the coral snake as a defense mechanism, with the key visual difference being their upturned snout (used for burrowing). The hognoses are, actually, mildly rear-fanged venomous; although generally not considered serious to humans, anecdotal research has suggested the South American hognoses to possess slightly more potent venom than the North American hognoses. In the Old World, none of the coral snake species usually fit the mnemonic.
Most species of coral snake are small in size. North American species average around in length, but specimens of up to or slightly larger have been reported. Some coral snakes even live in the water but most of them do not. Aquatic species have flattened tails that act as fins, aiding in swimming.
Behavior
Coral snakes vary widely in their behavior, but most are very elusive, fossorial (burrowing) snakes which spend most of their time buried beneath the ground or in the leaf litter of a rainforest floor, coming to the surface only when it rains or during breeding season. Some species, like Micrurus surinamensis, are almost entirely aquatic and spend most of their lives in slow-moving bodies of water that have dense vegetation.
Coral snakes feed mostly on smaller snakes, lizards, frogs, nestling birds, small rodents, etc.
Like all elapid snakes, coral snakes possess a pair of small hollow fangs to deliver their venom. The fangs are positioned at the front of the mouth. The fangs are fixed in position rather than retractable, and rather than being directly connected to the venom duct, they have a small groove through which the venom enters the base of the fangs. Because the fangs are relatively small and inefficient for venom delivery, rather than biting quickly and letting go (like vipers), coral snakes tend to hold onto their prey and make chewing motions when biting. The venom takes time to reach full effect.
Coral snakes are not aggressive or prone to biting and account for less than one percent of the total number of snake bites each year in the United States. The life span of coral snakes in captivity is about seven years.
Reproduction
M. fulvius reproduction is internal fertilization through the use of hemipenes. The breeding season occurs from spring to early summer and late summer to early fall. Male combat is not typical in M. fulvius as males are smaller than females. Micrurus fulvius are oviparous and typically lay eggs from May to July. During early spring females will undergo sudden vitellogenesis–oocyte and yolk formation–in preparation for breeding. Approximately 37 days post fertilization oviposition occurs and the average clutch size ranges from five to seven eggs. However, some in captivity have laid as many as thirteen eggs. The incubation period of the M. fulvius eggs normally reaches 60 days. Males also undergo sexual changes throughout the year, testicular recrudescence start in the fall and testicular regression occurs come spring. However, males typically have mature sperm residing in the epididymis year round and are capable of storing sperm in the deferent duct over the winter till the females are receptive. A study investigating how climate influences the reproductive cycle discovered species found closer to the equator displayed more continuous cycles while those in colder regions had more seasonal cycles. With increasing temperatures as a result of climate change, continuous cycles have the possibility of becoming more prevalent. Offspring reach maturation depending on sex, males mature at roughly 11 to 16 months while females reach maturity later at 26 months.
Distribution (U.S.)
New World coral snakes exist in the southern range of many temperate U.S. states. Coral snakes are found in scattered localities in the southern coastal plains from North Carolina to Louisiana, including all of Florida. They can be found in pine and scrub oak sandhill habitats in parts of this range, but sometimes inhabit hardwood areas and pine flatwoods that undergo seasonal flooding.
There is controversy about the classification of the very similar Texas coral snake as a separate species. Its habitat, in Texas, Louisiana, and Arkansas is separated from the eastern coral snake's habitat by the Mississippi River. The coral snake population is most dense in the southeastern United States.
The Arizona coral snake is classified as a separate species and genus and is found in central and southern Arizona, extreme southwestern New Mexico and southward to Sinaloa in western Mexico. It occupies arid and semiarid regions in many different habitat types, including thornscrub, desert-scrub, woodland, grassland and farmland. It is found in the plains and lower mountain slopes, at elevations ranging from sea level to ; often found in rocky areas.
Danger to humans
New World coral snakes possess one of the most potent venoms of any North American snake. However, relatively few bites are recorded due to their reclusive nature and the fact they generally inhabit sparsely populated areas. Even in areas that are densely populated, bites are rare. According to the American National Institutes of Health, there are an average of 15–25 coral snake bites in the United States each year.
When confronted by humans, coral snakes will almost always attempt to flee, and bite only if restrained. In addition, coral snakes have short fangs (proteroglyph dentition) that cannot penetrate thick clothing although bites are possible through normal thin clothing. Any skin penetration, however, is a medical emergency that requires immediate attention.
Historically, the venom of the North American Micrurus and Micruroides species was believed to contain powerful neurotoxins which could paralyze the breathing muscles, requiring mechanical or artificial respiration. It was usually reported that there was only mild pain associated with a bite and that respiratory failure could occur and onset of clinical symptoms may be delayed for as much as 10 to 18 hours. Coral Snake venom contains a neurotoxic component and symptoms include pain from the bite, sialorrhea, paresthesia, ptosis, weakness, blurred vision, paralysis, fasciculation and diplopia. However recent studies on the bites of the Texas coral snake (Micrurus tener) have shown that these bites rarely require antivenom, don't usually show any systemic respiratory problems and can be intensely painful. Further studies are necessary to see if these clinical features are true of all Micrurus species. Body length in coral snakes appears to be positively associated with venom yield.
Shortages of coral snake antivenom were previously reported, but one source states that production has resumed and, Pfizer indicates that antivenom is available.
Old World
Genus Calliophis
Species in this genus are:
Calliophis beddomei (M.A. Smith, 1943) – Beddome's coral snake (India)
Calliophis bibroni (Jan, 1858) – Bibron's coral snake (India)
Calliophis bilineatus
Calliophis bivirgatus (F. Boie, 1827) – blue Malaysian coral snake (Indonesia, Cambodia, Malaysia, Singapore, Thailand)
Calliophis castoe E.N. Smith, Ogale, Deepak & Giri, 2012 – Castoe's coral snake (India)
Calliophis gracilis (Gray, 1835) – spotted coral snake (Thailand, Malaysia, Indonesia, Singapore)
Calliophis haematoetron E.N. Smith, Manamendra-Arachchi & Somaweera, 2008 – blood-bellied coral snake (Sri Lanka)
Calliophis intestinalis (Laurenti, 1768) – banded Malaysian coral snake (Indonesia, Malaysia)
Calliophis maculiceps (Günther, 1858) – speckled coral snake (Myanmar, Thailand, Malaysia, Vietnam, Cambodia, Laos)
Calliophis melanurus (Shaw, 1802) – Indian coral snake (India, Bangladesh, Sri Lanka)
Calliophis nigrescens (Günther, 1862) – black coral snake (India)
Calliophis nigrotaeniatus
Calliophis philippinus
Calliophis salitan R.M. Brown, Smart, Leviton & E.N. Smith, 2018 – Dinagat Island Banded Coralsnake (Philippines)
Calliophis suluensis
Nota bene: A binomial authority in parentheses indicates that the species was originally described in a genus other than Calliophis.
Genus Hemibungarus
Species in this genus are:
Hemibungarus calligaster (Wiegmann, 1835) – barred coral snake (Philippines)
Hemibungarus gemianulis W. Peters, 1872 – (Philippines)
Hemibungarus mcclungi
Genus Sinomicrurus
Species in this genus are:
Sinomicrurus annularis
Sinomicrurus boettgeri
Sinomicrurus iwasakii
Sinomicrurus japonicus (Günther, 1868) – Japanese coral snake (Ryukyu Islands)
Sinomicrurus kelloggi (Pope, 1928) – Kellogg's coral snake (Vietnam, Laos, China)
Sinomicrurus macclellandi (J.T. Reinhardt, 1844) – Macclelland's coral snake (India, Nepal, Myanmar, Thailand, Vietnam, China, Ryukyu Islands, Taiwan)
Sinomicrurus peinani
Sinomicrurus sauteri (Steindachner, 1913) (Taiwan)
Sinomicrurus swinhoei
New World
Genus Micruroides
Micruroides euryxanthus (Kennicott, 1860) – Arizona coral snake (lowland regions from Arizona to Sinaloa, Mexico)
Micruroides euryxanthus australis Zweifel & Norris, 1955
Micruroides euryxanthus euryxanthus (Kennicott, 1860)
Micruroides euryxanthus neglectus Roze, 1967
Genus Micrurus
Nota bene: In the following list, a binomial authority or a trinomial authority in parentheses indicates that the species or subspecies was originally described in a genus other than Micrurus.
Micrurus albicinctus Amaral, 1925 – White-banded Coral Snake
Micrurus alleni K.P. Schmidt, 1936 – Allen's coral snake (eastern Nicaragua, Costa Rica, and Panama)
Micrurus alleni alleni K.P. Schmidt, 1936
Micrurus alleni richardi Taylor, 1951
Micrurus alleni yatesi Taylor, 1954
Micrurus altirostris (Cope, 1860) (Brazil, Uruguay, and northeastern Argentina)
Micrurus ancoralis Jan, 1872 – regal coral snake (southeastern Panama, western Colombia, and western Ecuador)
Micrurus ancoralis ancoralis Jan 1872
Micrurus ancoralis jani K.P. Schmidt, 1936
Micrurus annellatus (W. Peters, 1871) – annellated coral snake (southeastern Ecuador, eastern Peru, Bolivia, and western Brazil)
Micrurus annellatus annellatus (W. Peters, 1871)
Micrurus annellatus balzanii (Boulenger, 1898)
Micrurus annellatus bolivianus Roze, 1967
Micrurus averyi K.P. Schmidt, 1939 – black-headed coral snake
Micrurus baliocoryphus (Cope, 1860) – Mesopotamian coral snake
Micrurus bernadi (Cope, 1887) (Mexico)
Micrurus bocourti (Jan 1872) – Ecuadorian coral snake (western Ecuador to northern Colombia)
Micrurus bogerti Roze, 1967 – Bogert's coral snake (Oaxaca)
Micrurus boicora Bernarde, Turci, Abegg & Franco, 2018 – Boicora Coral Snake
Micrurus brasiliensis Roze, 1967 – Brazilian short-tailed coral snake
Micrurus browni K.P. Schmidt & H.M. Smith, 1943 – Brown's coral snake (Quintana Roo to Honduras)
Micrurus browni browni K.P. Schmidt & H.M. Smith, 1943
Micrurus browni importunus Roze, 1967
Micrurus browni taylori K.P. Schmidt & H.M. Smith, 1943
Micrurus camilae Renjifo & Lundberg, 2003 (Colombia)
Micrurus catamayensis Roze, 1989 – Catamayo coral snake (Catamayo Valley of Ecuador)
Micrurus clarki K.P. Schmidt, 1936 – Clark's coral snake (southeastern Costa Rica to western Colombia)
Micrurus collaris (Schlegel, 1837) – Guyana blackback coral snake (northern South America)
Micrurus collaris collaris (Schlegel, 1837)
Micrurus collaris breviventris Roze & Bernal-Carlo, 1987
Micrurus corallinus (Merrem, 1820) – painted coral snake
Micrurus decoratus (Jan 1858) – Brazilian coral snake
Micrurus diana Roze, 1983 – Diana's coral snake
Micrurus diastema (A.M.C. Duméril, Bibron & A.H.A. Duméril, 1854) – variable coral snake
Micrurus diastema aglaeope (Cope, 1859)
Micrurus diastema alienus (F. Werner, 1903)
Micrurus diastema affinis (Jan 1858)
Micrurus diastema apiatus (Jan 1858)
Micrurus diastema diastema (A.M.C. Duméril, Bibron & A.H.A. Duméril, 1854)
Micrurus diastema macdougalli Roze, 1967
Micrurus diastema sapperi (F. Werner, 1903)
Micrurus dissoleucus (Cope, 1860) – pygmy coral snake
Micrurus dissoleucus dissoleucus (Cope, 1860)
Micrurus dissoleucus dunni Barbour, 1923
Micrurus dissoleucus melanogenys (Cope, 1860)
Micrurus dissoleucus meridensis Roze, 1989
Micrurus dissoleucus nigrirostris K.P. Schmidt, 1955
Micrurus distans (Kennicott, 1860) – West Mexican coral snake
Micrurus distans distans (Kennicott, 1860)
Micrurus distans michoacanensis (Dugės, 1891)
Micrurus distans oliveri Roze, 1967
Micrurus distans zweifeli Roze, 1967
Micrurus diutius Burger, 1955 – Trinidad Ribbon Coral Snake
Micrurus dumerilii (Jan 1858)
Micrurus dumerilii antioquiensis K.P. Schmidt, 1936
Micrurus dumerilii carinicaudus K.P. Schmidt, 1936
Micrurus dumerilii colombianus (Griffin, 1916)
Micrurus dumerilii dumerilii (Jan 1858)
Micrurus dumerilii transandinus K.P. Schmidt, 1936
Micrurus dumerilii venezuelensis Roze, 1989
Micrurus elegans (Jan 1858) – elegant coral snake
Micrurus elegans elegans (Jan 1858)
Micrurus elegans veraepacis K.P. Schmidt, 1933
Micrurus ephippifer (Cope, 1886) – Oaxacan coral snake
Micrurus ephippifer ephippifer (Cope, 1886)
Micrurus ephippifer zapotecus Roze, 1989
Micrurus filiformis (Günther, 1859) – slender coral snake
Micrurus filiformis filiformis (Günther, 1859)
Micrurus filiformis subtilis Roze, 1967
Micrurus frontalis (A.M.C. Duméril, Bibron & A.H.A. Duméril, 1854) – southern coral snake (Brazil to northeastern Argentina)
Micrurus frontalis frontalis (A.M.C. Duméril, Bibron & A.H.A. Duméril, 1854)
Micrurus frontalis mesopotamicus Barrio & Miranda, 1967
Micrurus fulvius (Linnaeus, 1766) – eastern coral snake (U.S. coastal plains of North Carolina to Louisiana)
Micrurus hemprichii (Jan 1858) – Hemprich's coral snake/Orange-banded coral snake
Micrurus hemprichii hemprichii (Jan 1858)
Micrurus hemprichii ortoni K.P. Schmidt, 1953
Micrurus hemprichii rondonianus Roze & Silva, 1990
Micrurus hippocrepis (W. Peters, 1862) – Mayan coral snake
Micrurus ibiboboca (Merrem, 1820) – Caatinga coral snake
Micrurus isozonus (Cope, 1860) – Venezuela coral snake
Micrurus langsdorffi (Wagler, 1824) – Langsdorff's coral snake
Micrurus laticollaris (W. Peters, 1870) – Balsan coral snake
Micrurus laticollaris laticollaris (W. Peters, 1870)
Micrurus laticollaris maculirostris Roze, 1967
Micrurus latifasciatus K.P. Schmidt, 1933 – broad-ringed coral snake
Micrurus lemniscatus (Linnaeus, 1758) – South American coral snake (most of low-lying areas of South America)
Micrurus lemniscatus carvalhoi Roze, 1967
Micrurus lemniscatus frontifasciatus (F. Werner, 1927)
Micrurus lemniscatus helleri K.P. Schmidt & F.J.W. Schmidt, 1925
Micrurus lemniscatus lemniscatus (Linnaeus, 1758)
Micrurus limbatus Fraser, 1964 – Tuxtlan coral snake
Micrurus limbatus limbatus Fraser, 1964
Micrurus limbatus spilosomus Pérez-Higareda & H.M. Smith, 1990
Micrurus margaritiferus Roze, 1967 – Speckled coral snake
Micrurus medemi Roze, 1967 – Villavicencio coral snake
Micrurus meridensis (Roze, 1989) – Merida's coral snake
Micrurus mertensi K.P. Schmidt, 1936 – Merten's coral snake
Micrurus mipartitus (A.M.C. Duméril, Bibron & A.H.A. Duméril, 1854) – redtail coral snake
Micrurus mipartitus anomalus (Boulenger, 1896)
Micrurus mipartitus decussatus (A.M.C. Duméril, Bibron, & A.H.A. Duméril, 1854)
Micrurus mipartitus mipartitus (A.M.C. Duméril, Bibron & A.H.A. Duméril, 1854)
Micrurus mipartitus semipartitus (Jan 1858)
Micrurus mosquitensis Schmidt, 1933 – Misquito coral snake
Micrurus multifasciatus (Jan 1858) – Many-banded coral snake
Micrurus multifasciatus multifasciatus (Jan 1858)
Micrurus multifasciatus hertwigi (F. Werner, 1897)
Micrurus multiscutatus Rendahl & Vestergren, 1940 – Cauca coral snake
Micrurus narduccii (Jan, 1863) – Andean blackback coral snake
Micrurus narduccii narduccii (Jan 1863)
Micrurus narduccii melanotus (W. Peters, 1881)
Micrurus nattereri (Schmidt, 1952) – Natterer's Coral Snake
Micrurus nebularis Roze, 1989 – Cloud forest coral snake
Micrurus nigrocinctus (Girard, 1854) – Central American coral snake (Yucatan and Chiapas to Colombia as well as western Caribbean islands)
Micrurus nigrocinctus babaspul Roze, 1967
Micrurus nigrocinctus coibensis K.P. Schmidt, 1936
Micrurus nigrocinctus divaricatus (Hallowell, 1855)
Micrurus nigrocinctus nigrocinctus (Girard, 1854)
Micrurus nigrocinctus ovandoensis K.P. Schmidt & H.M. Smith, 1943
Micrurus nigrocinctus wagneri Mertens, 1941
Micrurus nigrocinctus yatesi Dunn, 1942
Micrurus nigrocinctus zunilensis K.P. Schmidt, 1932
Micrurus obscurus (Jan 1872) – Bolivian coral snake
Micrurus oligoanellatus (Ayerbe & Lopez, 2005) – Tambito's coral snake
Micrurus ornatissimus (Jan 1858) – Ornate Coral Snake
Micrurus pacaraimae Morato de Carvalho, 2002 – Pacaraima coral snake
Micrurus pachecogili Campbell, 2000 – Pueblan coral snake
Micrurus paraensis da Cunha & Nascimento, 1973 – Pará coral snake
Micrurus peruvianus K.P. Schmidt, 1936 – Peruvian coral snake
Micrurus petersi Roze, 1967 – Peters' coral snake
Micrurus potyguara Pires, Silva, Feitosa, Prudente, Pereira-Filho & Zaher, 2014 – Potyguara coral snake
Micrurus proximans H.M. Smith & Chrapliwy, 1958 – Nayarit coral snake
Micrurus psyches (Daudin, 1803) – Carib coral snake
Micrurus psyches circinalis (A.M.C. Duméril, Bibron & A.H.A. Duméril, 1854)
Micrurus psyches donosoi Hoge, Cordeiro & Romano, 1976
Micrurus psyches psyches (Daudin, 1803)
Micrurus putumayensis Lancini, 1962 – Putumayo coral snake
Micrurus pyrrhocryptus (Cope, 1862) – Argentinian coral snake (Argentina)
Micrurus remotus Roze, 1987 – Remote coral snake
Micrurus renjifoi Lamar, 2003 – Renjifo’s coral snake
Micrurus ruatanus (Günther, 1895) – Roatán coral snake
Micrurus sangilensis Nicéforo-María, 1942 – Santander coral snake
Micrurus scutiventris (Cope, 1869)
Micrurus serranus (Harvey, Aparicio & Gonzalez, 2003)
Micrurus silviae Di-Bernardo, Borges-Martins & Silva, 2007
Micrurus spixii (Wagler, 1824) – Amazon coral snake
Micrurus spixiii martiusi K.P. Schmidt, 1953
Micrurus spixii obscurus (Jan 1872)
Micrurus spixii princeps (Boulenger, 1905)
Micrurus spixii spixii (Wagler, 1824)
Micrurus spurelli (Boulenger, 1914)
Micrurus steindachneri (F. Werner, 1901) – Steindachner's coral snake
Micrurus steindachneri orcesi Roze, 1967
Micrurus steindachneri steindachneri (F. Werner, 1901)
Micrurus stewarti Barbour & Amaral, 1928 - Panamanian coral snake
Micrurus stuarti Roze, 1967 – Stuart's coral snake
Micrurus surinamensis (Cuvier, 1817) - Aquatic coral snake
Micrurus surinamensis nattereri K.P. Schmidt, 1952
Micrurus surinamensis surinamensis (Cuvier, 1817)
Micrurus tener (Baird & Girard, 1853) – Texas coral snake (Texas and Louisiana south to Morelos and Guanajuato)
Micrurus tener fitzingeri (Jan 1858)
Micrurus tener maculatus Roze, 1967
Micrurus tener microgalbineus B.C. Brown & H.M. Smith, 1942
Micrurus tener tamaulipensis Lavin-Murcio & Dixon, 2004
Micrurus tener tener (Baird & Girard, 1853)
Micrurus tikuna Feitosa, Silva, Pires, Zaher & Prudente, 2015
Micrurus tricolor
Micrurus tschudii (Jan 1858) – desert coral snake
Micrurus tschudii olssoni K.P. Schmidt & F.J.W. Schmidt, 1925
Micrurus tschudii tschudii (Jan 1858)
Mimicry
New World coral snakes serve as models for their Batesian mimics, false coral snakes, snake species whose venom is less toxic, as well as for many nonvenomous snake species that bear superficial resemblances to them. The role of coral snakes as models for Batesian mimics is supported by research showing that coral snake color patterns deter predators from attacking snake-shaped prey, and that in the absence of coral snakes, species hypothesized to mimic them are indeed attacked more frequently. Species that appear similar to coral snakes include:
Cemophora coccinea
Chionactis palarostris
Erythrolamprus aesculapii
Erythrolamprus bizona
Erythrolamprus ocellatus, Tobago false coral snake
Lampropeltis elapsoides, scarlet kingsnake
Lampropeltis pyromelana
Lampropeltis triangulum, milk snake, including the following subspecies and others:
Lampropeltis triangulum amaura
Lampropeltis triangulum annulata
Lampropeltis triangulum campbelli
Lampropeltis triangulum gaigeae
Lampropeltis triangulum gentilis
Lampropeltis triangulum hondurensis
Lampropeltis triangulum multistrata
Lampropeltis triangulum syspila
Lampropeltis zonata
Lystrophis pulcher, tri-color hognose snake
Oxyrhopus petola
Oxyrhopus rhombifer, false coral snake
Pliocercus elapoides, variegated false coral snake
Rhinobothryum bovallii, coral mimic snake, false tree coral
Rhinocheilus lecontei tessellatus
| Biology and health sciences | Reptiles | null |
611537 | https://en.wikipedia.org/wiki/Magnesite | Magnesite | Magnesite is a mineral with the chemical formula (magnesium carbonate). Iron, manganese, cobalt, and nickel may occur as admixtures, but only in small amounts.
Occurrence
Magnesite occurs as veins in and an alteration product of ultramafic rocks, serpentinite and other magnesium rich rock types in both contact and regional metamorphic terrains. These magnesites are often cryptocrystalline and contain silica in the form of opal or chert.
Magnesite is also present within the regolith above ultramafic rocks as a secondary carbonate within soil and subsoil, where it is deposited as a consequence of dissolution of magnesium-bearing minerals by carbon dioxide in groundwaters.
Formation
Magnesite can be formed via talc carbonate metasomatism of peridotite and other ultramafic rocks. Magnesite is formed via carbonation of olivine in the presence of water and carbon dioxide at elevated temperatures and high pressures typical of the greenschist facies.
Magnesite can also be formed via the carbonation of magnesium serpentine (lizardite) via the following reaction:
2 Mg3Si2O5(OH)4 + 3 CO2 → Mg3Si4O10(OH)2 + 3 MgCO3 + 3 H2O
However, when performing this reaction in the laboratory, the trihydrated form of magnesium carbonate (nesquehonite) will form at room temperature. This very observation led to the postulation of a "dehydration barrier" being involved in the low-temperature formation of anhydrous magnesium carbonate. Laboratory experiments with formamide, a liquid resembling water, have shown how no such dehydration barrier can be involved. The fundamental difficulty to nucleate anhydrous magnesium carbonate remains when using this non-aqueous solution. Not cation dehydration, but rather the spatial configuration of carbonate anions creates the barrier in the low-temperature nucleation of magnesite.
Magnesite has been found in modern sediments, caves and soils. Its low-temperature (around ) formation is known to require alternations between precipitation and dissolution intervals. The low-temperature formation of magnesite might well be of significance toward large-scale carbon sequestration. A major step forward toward the industrial production of magnesite at atmospheric pressure and a temperature of 316 K was described by Vandeginste. In those experiments small additions of hydrochloric acid alternated periodically with additions of sodium carbonate solution. New was also the very short duration of only a few hours for the alternating dissolution and precipitation cycles.
Magnesite was detected in meteorite ALH84001 and on planet Mars itself. Magnesite was identified on Mars using infrared spectroscopy from satellite orbit. Near Jezero Crater, Mg-carbonates have been detected and reported to have formed in lacustrine environment prevailing there. Controversy still exists over the temperature of formation of these carbonates. Low-temperature formation has been suggested for the magnesite from the Mars-derived ALH84001 meteorite.
Magnesium-rich olivine (forsterite) favors production of magnesite from peridotite. Iron-rich olivine (fayalite) favors production of magnetite-magnesite-silica compositions.
Magnesite can also be formed by way of metasomatism in skarn deposits, in dolomitic limestones, associated with wollastonite, periclase, and talc.
Resistant to high temperature and able to withstand high pressure, magnesite has been proposed to be one of the major carbonate bearing phase in Earth's mantle and possible carriers for deep carbon reservoirs. For similar reason, it is found in metamorphosed peridotite rocks in Central Alps, Switzerland and high pressure eclogitic rocks from Tianshan, China.
Magnesite can also precipitate in lakes in presence of bacteria either as hydrous Mg-carbonates or magnesite.
Uses
Refractory material
Similar to the production of lime, magnesite can be burned in the presence of charcoal to produce MgO, which, in the form of a mineral, is known as periclase. Large quantities of magnesite are burnt to make magnesium oxide: an important refractory (heat-resistant) material used as a lining in blast furnaces, kilns and incinerators.
Calcination temperatures determine the reactivity of resulting oxide products and the classifications of light burnt and dead burnt refer to the surface area and resulting reactivity of the product (this is typically determined by an industry metric of the iodine number).
'Light burnt' product generally refers to calcination commencing at 450 °C and proceeding to an upper limit of 900 °C – which results in good surface area and reactivity.
Above 900 °C, the material loses its reactive crystalline structure and reverts to the chemically inert 'dead-burnt' product- which is preferred for use in refractory materials such as furnace linings.
In fire assay, magnesite cupels can be used for cupellation, as the magnesite cupel will resist the high temperatures involved.
Other uses
Magnesite can also be used as a binder in flooring material (magnesite screed). Furthermore, it is being used as a catalyst and filler in the production of synthetic rubber and in the preparation of magnesium chemicals and fertilizers.
Research is proceeding to evaluate the practicality of sequestering the greenhouse gas carbon dioxide in magnesite on a large scale. This has focused on peridotites from ophiolites (obducted mantle rocks on crust) where magnesite can be created by letting carbon dioxide react with these rocks. Some progress has been made in ophiolites from Oman. But the major problem is that these artificial processes require sufficient porosity-permeability so that the fluids can flow but this is hardly the case in peridotites.
Artworks
Magnesite can be cut, drilled, and polished to form beads that are used in jewelry-making. Magnesite beads can be dyed into a broad spectrum of bold colors, including a light blue color that mimics the appearance of turquoise.
The Japanese-American artist Isamu Noguchi used magnesite as a sculptural material for some of his artworks.
Isotopic structure
The recent advancement in the field of stable isotope geochemistry is the study of isotopic structure of minerals and molecules. This requires study of molecules with high resolutions looking at bonding scenario (how heavy isotopes are bonded to each other)- leading to knowledge of stability of molecule depending on its isotopic structure.
Isotopically substituted molecules have higher mass. As a consequence, molecular vibration reduces and the molecule develops a lower zero point energy (see Kinetic isotope effect).
The abundances of certain bonds in certain molecules are sensitive to temperature at which it formed (e.g., abundance of 13C16O18O in carbonates as 13C-18O bond). This information has been exploited to form the foundation of clumped isotope geochemistry. Clumped isotope thermometers have been established for carbonate minerals like dolomite, calcite, siderite etc and non-carbonate compounds like methane and oxygen. Depending on the strength of cation-carbonate oxygen (ie, Mg-O, Ca-O) bonds- different carbonate minerals can form or preserve clumped isotopic signatures differently.
Measurements and reporting
Clumped isotopic analysis has certain aspects to it. These are:
Digestion, analysis and acid fractionation correction
Clumped isotopic analysis is usually done by gas source mass spectrometry where the CO2 liberated from magnesite by phosphoric acid digestion is fed into the isotope ratio mass spectrometer. In such scenario, one needs to ensure that liberation of CO2 from magnesite is complete. Digesting magnesite is hard since it takes a long time and different labs report different digestion times and temperatures (from 12 hours at 100 °C to 1 hour at 90 °C in phosphoric acid). Due to digestion at this high temperature, some of the 13C-18O bonds in the liberated CO2 are broken (leading to reduction in abundance of 'clumped' CO2) during phosphoric acid digestion of carbonates. To account for this additional (analytical artifact), a correction called the 'acid fractionation correction' is added to the magnesite clumped isotope value obtained at temperature of digestion.
Since the CO2 gas is liberated from carbonate mineral during acid digestion, leaving one O behind- a fractionation occurs, and the isotopic composition of the analyzed CO2 gas needs to be corrected for this. For magnesite, the most reliable fractionation factor(α) equation is given as:
103ln(α) = [(6.845 ± 0.475)∗105/T2] + (4.22 ± 0.08); T in K
Different researchers have also used other fractionation factors like dolomite fractionation factor.
Standards
While measuring samples of unknown composition, it is required to measure some standard materials (see Reference materials for stable isotope analysis). With internal standards and reference materials, analytical session is routinely monitored. Standard materials are majorly calcite and marble.
Δ47 – Temperature calibration
To convert clumped isotope data into temperature, a calibration curve is required which expresses the functional form of temperature dependence of clumped isotope composition. No mineral specific calibration exists for magnesite. Based on some experimental data where mineral precipitation temperature and clumped isotope derived temperature doesn't match, a need of mineral specific calibration emerges. The mismatch arises since bonding in magnesite is different from calcite/dolomite and/or acid digestion is conducted at higher temperature.
Magnesite-water and CO2-magnesite isotope fractionation factors
Using clumped isotope derived temperature, C and O isotopic composition of the parental fluid can be calculated using known magnesite-fluid isotope fractionation factors, since fractionation is temperature dependent. Reported magnesite-fluid O and C isotope fractionation factors in literature are not in agreement with each other. The fractionation behaviors have not been substantiated by experimental observation.
Factors controlling isotopic structure in magnesite
Conversion from hydrous Mg-carbonates to magnesite
In low temperature, thus, hydrous Mg-carbonates (hydromagnesite, nesquehonite etc.) form. It is possible to convert these phases into magnesite by changing temperature by mineral dissolution-precipitation or dehydration. While so happens, an isotope effect associated can control the isotopic composition of precipitated magnesite.
Disequilibrium
Disequilibrium processes like degassing, rapid CO2 uptake etc. modify clumped isotopic composition of carbonate minerals specifically at low temperatures. They variably enrich or deplete the system in heavy isotopes of C and O. Since clumped isotope abundance depends on abundance of isotopes of C and O, they are also modified. Another very prominent effect here is that of pH of precipitating fluid. As pH of precipitating fluid changes, DIC pool is affected and isotopic composition of precipitating carbonate changes.
Mineral structure and later thermal effects
Crystalline and cryptocrystalline magnesites have very different mineral structures. While crystalline magnesite has a well developed crystal structure, the cryptocrystalline magnesite is amorphous- mostly aggregate of fine grains. Since clumped isotopic composition depends on specific bonding, difference in crystal structure is very likely to affect the way clumped isotopic signatures are recorded in these different structures. This leads to the fact that their pristine signatures might be modified differently by later thermal events like diagenesis/burial heating etc.
Information on formation from isotopic structure
Clumped isotopes have been used in interpreting conditions of magnesite formation and the isotopic composition of the precipitating fluid. Within ultramafic complexes, magnesites are found within veins and stockworks in cryptocrystalline form as well as within carbonated peridotite units in crystalline form. These cryptocrystalline forms are mostly variably weathered and yield low temperature of formation. On the other hand, coarse magnesites yield very high temperature indicating hydrothermal origin. It is speculated that coarse high temperature magnesites are formed from mantle derived fluids whereas cryptocrystalline ones are precipitated by circulating meteoric water, taking up carbon from dissolved inorganic carbon pool, soil carbon and affected by disequilibrium isotope effects.
Magnesites forming in lakes and playa settings are in general enriched in heavy isotopes of C and O because of evaporation and CO2 degassing. This reflects in the clumped isotope derived temperature being very low. These are affected by pH effect, biological activity as well as kinetic isotope effect associated with degassing. Magnesite forms as surface moulds in such conditions but more generally occur as hydrous Mg-carbonates since their precipitation is kinetically favored. Most of the times, they derive C from DIC or nearby ultramafic complexes (e.g., Altin Playa, British Columbia, Canada).
Magnesites in metamorphic rocks, on the other hand, indicate very high temperature of formation. Isotopic composition of parental fluid is also heavy- generally metamorphic fluids. This has been verified by fluid inclusion derived temperature as well as traditional O isotope thermometry involving co-precipitating quartz-magnesite.
Often, magnesite records lower clumped isotope temperature than associated dolomite, calcite. The reason might be that calcite, dolomite form earlier at higher temperature (from mantle like fluids) which increases Mg/Ca ratio in the fluid sufficiently so as to precipitate magnesite. As this happens with increasing time, fluid cools, evolves by mixing with other fluids and when it forms magnesite, it decreases its temperature. So the presence of associated carbonates have a control on magnesite isotopic composition.
Origin of Martian carbonates can be deconvolved with the application of clumped isotope. Source of the CO2, climatic-hydrologic conditions on Mars could be assessed from these rocks. Recent study has shown (implementing clumped isotope thermometry) that carbonates in ALH84001 indicate formation at low temperature evaporative condition from subsurface water and derivation of CO2 from Martian atmosphere.
Occupational safety and health
People can be exposed to magnesite in the workplace by inhaling it, skin contact, and eye contact.
United States
The Occupational Safety and Health Administration (OSHA) has set the legal limit (permissible exposure limit) for magnesite exposure in the workplace as 15 mg/m3 total exposure and 5 mg/m3 respiratory exposure over an 8-hour workday. The National Institute for Occupational Safety and Health (NIOSH) has set a recommended exposure limit (REL) of 10 mg/m3 total exposure and 5 mg/m3 respiratory exposure over an 8-hour workday.
| Physical sciences | Minerals | Earth science |
611589 | https://en.wikipedia.org/wiki/Java%20Man | Java Man | Java Man (Homo erectus erectus, formerly also Anthropopithecus erectus or Pithecanthropus erectus) is an early human fossil discovered in 1891 and 1892 on the island of Java (Indonesia). Estimated to be between 700,000 and 1,490,000 years old, it was, at the time of its discovery, the oldest hominid fossil ever found, and it remains the type specimen for Homo erectus.
Led by Eugène Dubois, the excavation team uncovered a tooth, a skullcap, and a thighbone at Trinil on the banks of the Solo River in East Java. Arguing that the fossils represented the "missing link" between apes and humans, Dubois gave the species the scientific name Anthropopithecus erectus, then later renamed it Pithecanthropus erectus. The fossil aroused much controversy. Within a decade of the discovery almost eighty books or articles had been published on Dubois's finds. Despite Dubois's argument, few accepted that Java Man was a transitional form between apes and humans. Some dismissed the fossils as apes and others as modern humans, whereas many scientists considered Java Man as a primitive side branch of evolution not related to modern humans at all. In the 1930s Dubois made the claim that Pithecanthropus was built like a "giant gibbon", a much misinterpreted attempt by Dubois to prove that it was the "missing link". Eventually, similarities between Java Man and Sinanthropus pekinensis (Peking Man) led Ernst Mayr to rename both Homo erectus in 1950, placing them directly in the human evolutionary tree.
To distinguish Java Man from other Homo erectus populations, some scientists began to regard it as a subspecies, Homo erectus erectus, in the 1970s. Other fossils found in the first half of the twentieth century in Java at Sangiran and Mojokerto, all older than those found by Dubois, are also considered part of the species Homo erectus. The fossils of Java Man have been housed at the Rijksmuseum van Geologie en Mineralogie and later Naturalis in the Netherlands since 1900.
History of discoveries
Background
Charles Darwin had argued that humanity evolved in Africa, because this is where great apes like gorillas and chimpanzees lived. Though Darwin's claims have since been vindicated by the fossil record, they were proposed without any fossil evidence. Other scientific authorities disagreed with him, like Charles Lyell, a geologist, and Alfred Russel Wallace, who thought of a similar theory of evolution around the same time as Darwin. Because both Lyell and Wallace believed that humans were more closely related to gibbons or another great ape (the orangutans), they identified Southeast Asia as the cradle of humanity because this is where these apes lived. Dutch anatomist Eugène Dubois favored the latter theory, and sought to confirm it.
Trinil fossils
In October 1887, Dubois abandoned his academic career and left for the Dutch East Indies (present-day Indonesia) to look for the fossilized ancestor of modern man. Having received no funding from the Dutch government for his eccentric endeavorsince no one at the time had ever found an early human fossil while looking for ithe joined the Dutch East Indies Army as a military surgeon. Because of his work duties, it was only in July 1888 that he began to excavate caves in Sumatra. Having quickly found abundant fossils of large mammals, Dubois was relieved of his military duties (March 1889), and the colonial government assigned two engineers and fifty convicts to help him with his excavations. After he failed to find the fossils he was looking for on Sumatra, he moved on to Java in 1890.
Again assisted by convict laborers and two army corporals, Dubois began searching along the Solo River near Trinil in August 1891. His team soon excavated a molar (Trinil 1) and a skullcap (Trinil 2). Its characteristics were a long cranium with a sagittal keel and heavy browridge. Dubois first gave them the name Anthropopithecus ("man-ape"), as the chimpanzee was sometimes known at the time. He chose this name because a similar tooth found in the Siwalik Hills in India in 1878 had been named Anthropopithecus, and because Dubois first assessed the cranium to have been about , closer to apes than to humans.
In August 1892, a year later, Dubois's team found a long femur (thighbone) shaped like a human one, suggesting that its owner had stood upright. The femur bone was found 50 feet (approx. 15 meters) from the original find one year earlier. Believing that the three fossils belonged to a single individual, "probably a very aged female", Dubois renamed the specimen Anthropopithecus erectus. Only in late 1892, when he determined that the cranium measured about , did Dubois consider that his specimen was a transitional form between apes and humans. In 1894, he thus renamed it Pithecanthropus erectus ("upright ape-man"), borrowing the genus name Pithecanthropus from Ernst Haeckel, who had coined it a few years earlier to refer to a supposed "missing link" between apes and humans. This specimen has also been known as Pithecanthropus 1.
Comparisons with Peking Man
In 1927, Canadian Davidson Black identified two fossilized teeth he had found in Zhoukoudian near Beijing as belonging to an ancient human, and named his specimen Sinanthropus pekinensis, now better known as Peking Man. In December 1929, the first of several skullcaps was found on the same site, and it appeared similar but slightly larger than Java Man. Franz Weidenreich, who replaced Black in China after the latter's death in 1933, argued that Sinanthropus was also a transitional fossil between apes and humans, and was in fact so similar to Java's Pithecanthropus that they should both belong to the family Hominidae. Eugène Dubois categorically refused to entertain this possibility, dismissing Peking Man as a kind of Neanderthal, closer to humans than the Pithecanthropus, and insisting that Pithecanthropus belonged to its own superfamily, the Pithecanthropoidea.
Other discoveries on Java
After the discovery of Java Man, Berlin-born paleontologist G. H. R. von Koenigswald recovered several other early human fossils in Java. Between 1931 and 1933 von Koenigswald discovered fossils of Solo Man from sites along the Bengawan Solo River on Java, including several skullcaps and cranial fragments. In 1936, von Koenigswald discovered a juvenile skullcap known as the Mojokerto child in East Java. Considering the Mojokerto child skull cap to be closely related to humans, von Koenigswald wanted to name it Pithecanthropus modjokertensis (after Dubois's specimen), but Dubois protested that Pithecanthropus was not a human but an "ape-man".
Von Koenigswald also made several discoveries in Sangiran, Central Java, where more fossils of early humans were discovered between 1936 and 1941. Among the discoveries was a skullcap of similar size to that found by Dubois at the Trinil 2 site. Von Koenigswald's discoveries in Sangiran convinced him that all these skulls belonged to early humans. Dubois again refused to acknowledge the similarity. Ralph von Koenigswald and Franz Weidenreich compared the fossils from Java and Zhoukoudian and concluded that Java Man and Peking Man were closely related. Dubois died in 1940, still refusing to recognize their conclusion, and official reports remain critical of the Sangiran site's poor presentation and interpretation.
Early interpretations
More than 50 years after Dubois's find, Ralph von Koenigswald recollected that, "No other paleontological discovery has created such a sensation and led to such a variety of conflicting scientific opinions." The Pithecanthropus fossils were so immediately controversial that by the end of the 1890s, almost 80 publications had already discussed them.
Until the Taung childthe 2.8 million-year-old remains of an Australopithecus africanuswere discovered in South Africa in 1924, Dubois's and Koenigswald's discoveries were the oldest hominid remains ever found. Some scientists of the day suggested that Dubois's Java Man was a potential intermediate form between modern humans and the common ancestor we share with the other great apes. The current consensus of anthropologists is that the direct ancestors of modern humans were African populations of Homo erectus (Homo ergaster), rather than the Asian populations of the same species exemplified by Java Man and Peking Man.
Missing link theory
Dubois first published his find in 1894. Dubois's central claim was that Pithecanthropus was a transitional form between apes and humans, a so-called "missing link". Many disagreed. Some critics claimed that the bones were those of an upright walking ape, or that they belonged to a primitive human. This judgment made sense at a time when an evolutionary view of humanity had not yet been widely accepted, and scientists tended to view hominid fossils as racial variants of modern humans rather than as ancestral forms. After Dubois let a number of scientists examine the fossils in a series of conferences held in Europe in the 1890s, they started to agree that Java Man may be a transitional form after all, but most of them thought of it as "an extinct side branch" of the human tree that had indeed descended from apes, but not evolved into humans. This interpretation eventually imposed itself and remained dominant until the 1940s.
Dubois was bitter about this and locked the fossil up in a trunk until 1923 when he showed it to Ales Hrdlicka from the Smithsonian Institution. In response to critics who refused to accept that Java Man was a "missing link", in 1932 Dubois published a paper arguing that the Trinil bones looked like those of a "giant gibbon". Dubois's use of the phrase has been widely misinterpreted as a retraction, but it was intended an argument to support his claim that Pithecanthropus was a transitional form. According to Dubois, evolution occurred by leaps, and the ancestors of humanity had doubled their brain-to-body ratio on each leap. To prove that Java Man was the "missing link" between apes and humans, he therefore had to show that its brain-to-body ratio was double that of apes and half that of humans. The problem was that Java Man's cranial capacity was 900 cubic centimeters, about two-thirds of modern humans'.
Like many scientists who believed that modern humans evolved "Out of Asia", Dubois thought that gibbons were closest to humans among the great apes. To preserve the proportions predicted by his theory of brain evolution, Dubois argued that Java Man was shaped more like a gibbon than a human. Imagined "with longer arms and a greatly expanded chest and upper body", the Trinil creature became a gigantic ape of about , but "double cephalization of the anthropoid apes in general and half that of man". It was therefore halfway on the path to becoming a modern human. As Dubois concluded his 1932 paper: "I still believe, now more firmly than ever, that the Pithecanthropus of Trinil is the real 'missing link.'"
Reclassification as Homo erectus
Based on Weidenreich's work and on his suggestion that Pithecanthropus erectus and Sinanthropus pekinensis were connected through a series of interbreeding populations, German biologist Ernst Mayr reclassified them both as being part of the same species: Homo erectus. Mayr presented his conclusion at the Cold Spring Harbor Symposium in 1950, and this resulted in Dubois's erectus species being reclassified under the genus Homo. As part of the reclassification, Mayr included not only Sinanthropus and Pithecanthropus, but also Plesianthropus, Paranthropus, Javanthropus, and several other genera as synonyms, arguing that all human ancestors were part of a single genus (Homo), and that "never one more than one species of man existed on the earth at any one time". A "revolution in taxonomy", Mayr's single-species approach to human evolution was quickly accepted. It shaped paleoanthropology in the 1950s and lasted into the 1970s, when the African genus Australopithecus was accepted into the human evolutionary tree.
In the 1970s a tendency developed to regard the Javanese variety of H. erectus as a subspecies, Homo erectus erectus, with the Chinese variety being referred to as Homo erectus pekinensis.
Post-discovery analysis
Date of the fossils
Dubois's complete collection of fossils were transferred between 1895 and 1900 to what is now known as Naturalis, in Leiden in the Netherlands. The main fossil of Java Man, the skullcap cataloged as "Trinil 2", has been dated biostratigraphically, that is, by correlating it with a group of fossilized animals (a "faunal assemblage") found nearby on the same geological horizon, which is itself compared with assemblages from other layers and classified chronologically. Ralph von Koenigswald first assigned Java Man to the Trinil Fauna, a faunal assemblage that he composed from several Javanese sites. He concluded that the skullcap was about 700,000 years old, thus dating from the beginning of the Middle Pleistocene.
Though this view is still widely accepted, in the 1980s a group of Dutch paleontologists used Dubois's collection of more than 20,000 animal fossils to reassess the date of the layer in which Java Man was found. Using only fossils from Trinil, they called that new faunal assemblage the Trinil H. K. Fauna, in which H. K. stands for Haupt Knochenschicht, or "main fossil-bearing layer". This assessment dates the fossils of Java Man to between 900,000 and 1,000,000 years old. On the other hand, work published in 2014 gives a "maximum age of 0.54 ± 0.10 million years and a minimum age of 0.43 ± 0.05 million years" for Ar-Ar and luminescence dating of sediment in human-predated shell material from Trinil. Work continues on assessing the dating of this complex site.
Other fossils attest to the even earlier presence of H. erectus in Java. Sangiran 2 (named after its discovery site) may be as old as 1.66 Ma (million years). The controversial Mojokerto child, which Carl C. Swisher and Garniss Curtis once dated to 1.81 ± 0.04 Ma, has now been convincingly re-dated to a maximum age of 1.49 ± 0.13 Ma, that is, 1.49 million years with a margin of error of plus or minus 130,000 years.
Type specimen
The fossils found in Java are considered the type specimen for H. erectus. Because the fossils of Java Man were found "scattered in an alluvial deposit"they had been laid there by the flow of a riverdetractors doubted that they belonged to the same species, let alone the same individual. German pathologist Rudolf Virchow, for instance, argued in 1895 that the femur was that of a gibbon. Dubois had difficulty convincing his critics, because he had not attended the excavation, and could not explain specifically enough the exact location of the bones. Because the Trinil thighbone looks very much like that of a modern human, it might have been a "reworked fossil", that is, a relatively young fossil that was deposited into an older layer after its own layer had been eroded. For this reason, there is still dissent about whether all the Trinil fossils represent the same species.
Physical characteristics
Java Man was about tall and his thighbones show that he walked erect like modern humans. The femur is thicker than that of a modern human, indicating he was engaging in a lot of running. The skull was characterized by thick bones and a retreating forehead. The large teeth made the jaw large and jutting, with the lower lips overhanging the lower margin of the mandible, giving the impression of no chin. The browridges were straight and massive. At 900 cm3, his cranial capacity was smaller than that of later H. erectus specimens. However, he had humanlike teeth with large canines.
Judging from anatomical and archeological aspects as well as Java Man's ecological role, meat from vertebrates was likely an important part of their diet. Java Man, like other Homo erectus, was probably a rare species. There is evidence that Java Man used shell tools to cut meat. Java Man's dispersal through Southeast Asia coincides with the extirpation of the giant turtle Megalochelys, possibly due to overhunting as the turtle would have been an easy, slow-moving target which could have been stored for quite some time.
Material culture
H. erectus arrived in Eurasia approximately 1.8 million years ago, in an event considered to be the first African exodus. There is evidence that the Java population of H. erectus lived in an ever-wet forest habitat. More specifically the environment resembled a savannah, but was likely regularly inundated ("hydromorphic savanna"). The plants found at the Trinil excavation site included grass (Poaceae), ferns, Ficus, and Indigofera, which are typical of lowland rainforest.
Control of fire
The control of fire by Homo erectus is generally accepted by archaeologists to have begun some 400,000 years ago, with claims regarding earlier evidence finding increasing scientific support. Burned wood has been found in layers that carried the Java Man fossils in Trinil, dating to around from 500,000 to 830,000 BP. However, because Central Java is a volcanic region, the charring may have resulted from natural fires, and there is no conclusive proof that Homo erectus in Java controlled fire. It has been proposed that Java Man was aware of the use of fire, and that the frequent presence of natural fires may have allowed Java Man "opportunistic use [... that] did not create an archeologically visible pattern".
| Biology and health sciences | Homo | Biology |
612029 | https://en.wikipedia.org/wiki/Cyclic%20model | Cyclic model | A cyclic model (or oscillating model) is any of several cosmological models in which the universe follows infinite, or indefinite, self-sustaining cycles. For example, the oscillating universe theory briefly considered by Albert Einstein in 1930 theorized a universe following an eternal series of oscillations, each beginning with a Big Bang and ending with a Big Crunch; in the interim, the universe would expand for a period of time before the gravitational attraction of matter causes it to collapse back in and undergo a bounce.
Overview
In the 1920s, theoretical physicists, most notably Albert Einstein, considered the possibility of a cyclic model for the universe as an (everlasting) alternative to the model of an expanding universe. In 1922, Alexander Friedmann introduced the Oscillating Universe Theory. However, work by Richard C. Tolman in 1934 showed that these early attempts failed because of the cyclic problem: according to the second law of thermodynamics, entropy can only increase. This implies that successive cycles grow longer and larger. Extrapolating back in time, cycles before the present one become shorter and smaller culminating again in a Big Bang and thus not replacing it. This puzzling situation remained for many decades until the early 21st century when the recently discovered dark energy component provided new hope for a consistent cyclic cosmology. In 2011, a five-year survey of 200,000 galaxies and spanning 7 billion years of cosmic time confirmed that "dark energy is driving our universe apart at accelerating speeds."
One new cyclic model is the brane cosmology model of the creation of the universe, derived from the earlier ekpyrotic model. It was proposed in 2001 by Paul Steinhardt of Princeton University and Neil Turok of Cambridge University. The theory describes a universe exploding into existence not just once, but repeatedly over time. The theory could potentially explain why a repulsive form of energy known as the cosmological constant, which is accelerating the expansion of the universe, is several orders of magnitude smaller than predicted by the standard Big Bang model.
A different cyclic model relying on the notion of phantom energy was proposed in 2007 by Lauris Baum and Paul Frampton of the University of North Carolina at Chapel Hill.
Other cyclic models include conformal cyclic cosmology and loop quantum cosmology.
The Steinhardt–Turok model
In this cyclic model, two parallel orbifold planes or M-branes collide periodically in a higher-dimensional space. The visible four-dimensional universe lies on one of these branes. The collisions correspond to a reversal from contraction to expansion, or a Big Crunch followed immediately by a Big Bang. The matter and radiation we see today were generated during the most recent collision in a pattern dictated by quantum fluctuations created before the branes. After billions of years the universe reached the state we observe today; after additional billions of years it will ultimately begin to contract again. Dark energy corresponds to a force between the branes, and serves the crucial role of solving the monopole, horizon, and flatness problems. Moreover, the cycles can continue indefinitely into the past and the future, and the solution is an attractor, so it can provide a complete history of the universe.
As Richard C. Tolman showed, the earlier cyclic model failed because the universe would undergo inevitable thermodynamic heat death. However, the newer cyclic model evades this by having a net expansion each cycle, preventing entropy from building up. However, there remain major open issues in the model. Foremost among them is that colliding branes are not understood by string theorists, and nobody knows if the scale invariant spectrum will be destroyed by the big crunch. Moreover, as with cosmic inflation, while the general character of the forces (in the ekpyrotic scenario, a force between branes) required to create the vacuum fluctuations is known, there is no candidate from particle physics.
The Baum–Frampton model
This more recent cyclic model of 2007 assumes an exotic form of dark energy called phantom energy, which possesses negative kinetic energy and would usually cause the universe to end in a Big Rip. This condition is achieved if the universe is dominated by dark energy with a cosmological equation of state parameter satisfying the condition , for energy density and pressure p. By contrast, Steinhardt–Turok assume . In the Baum–Frampton model, a septillionth (or less) of a second (i.e. 10−24 seconds or less) before the would-be Big Rip, a turnaround occurs and only one causal patch is retained as our universe. The generic patch contains no quark, lepton or force carrier; only dark energy – and its entropy thereby vanishes. The adiabatic process of contraction of this much smaller universe takes place with constant vanishing entropy and with no matter including no black holes which disintegrated before turnaround.
The idea that the universe "comes back empty" is a central new idea of this cyclic model, and avoids many difficulties confronting matter in a contracting phase such as excessive structure formation, proliferation and expansion of black holes, as well as going through phase transitions such as those of QCD and electroweak symmetry restoration. Any of these would tend strongly to produce an unwanted premature bounce, simply to avoid violation of the second law of thermodynamics. The condition of may be logically inevitable in a truly infinitely cyclic cosmology because of the entropy problem. Nevertheless, many technical back up calculations are necessary to confirm consistency of the approach. Although the model borrows ideas from string theory, it is not necessarily committed to strings, or to higher dimensions, yet such speculative devices may provide the most expeditious methods to investigate the internal consistency. The value of in the Baum–Frampton model can be made arbitrarily close to, but must be less than, −1.
Other cyclic models
Conformal cyclic cosmology—a general relativity based theory by Roger Penrose in which the universe expands until all the matter decays and is turned to light—so there is nothing in the universe that has any time or distance scale associated with it. This permits it to become identical with the Big Bang, so starting the next cycle.
Loop quantum cosmology which predicts a "quantum bridge" between contracting and expanding cosmological branches.
| Physical sciences | Physical cosmology | Astronomy |
612057 | https://en.wikipedia.org/wiki/Potential%20well | Potential well | A potential well is the region surrounding a local minimum of potential energy. Energy captured in a potential well is unable to convert to another type of energy (kinetic energy in the case of a gravitational potential well) because it is captured in the local minimum of a potential well. Therefore, a body may not proceed to the global minimum of potential energy, as it would naturally tend to do due to entropy.
Overview
Energy may be released from a potential well if sufficient energy is added to the system such that the local maximum is surmounted. In quantum physics, potential energy may escape a potential well without added energy due to the probabilistic characteristics of quantum particles; in these cases a particle may be imagined to tunnel through the walls of a potential well.
The graph of a 2D potential energy function is a potential energy surface that can be imagined as the Earth's surface in a landscape of hills and valleys. Then a potential well would be a valley surrounded on all sides with higher terrain, which thus could be filled with water (e.g., be a lake) without any water flowing away toward another, lower minimum (e.g. sea level).
In the case of gravity, the region around a mass is a gravitational potential well, unless the density of the mass is so low that tidal forces from other masses are greater than the gravity of the body itself.
A potential hill is the opposite of a potential well, and is the region surrounding a local maximum.
Quantum confinement
Quantum confinement can be observed once the diameter of a material is of the same magnitude as the de Broglie wavelength of the electron wave function. When materials are this small, their electronic and optical properties deviate substantially from those of bulk materials.
A particle behaves as if it were free when the confining dimension is large compared to the wavelength of the particle. During this state, the bandgap remains at its original energy due to a continuous energy state. However, as the confining dimension decreases and reaches a certain limit, typically in nanoscale, the energy spectrum becomes discrete. As a result, the bandgap becomes size-dependent. As the size of the particles decreases, the electrons and electron holes come closer, and the energy required to activate them increases, which ultimately results in a blueshift in light emission.
Specifically, the effect describes the phenomenon resulting from electrons and electron holes being squeezed into a dimension that approaches a critical quantum measurement, called the exciton Bohr radius. In current application, a quantum dot such as a small sphere confines in three dimensions, a quantum wire confines in two dimensions, and a quantum well confines only in one dimension. These are also known as zero-, one- and two-dimensional potential wells, respectively. In these cases they refer to the number of dimensions in which a confined particle can act as a free carrier. See external links, below, for application examples in biotechnology and solar cell technology.
Quantum mechanics view
The electronic and optical properties of materials are affected by size and shape. Well-established technical achievements including quantum dots were derived from size manipulation and investigation for their theoretical corroboration on quantum confinement effect. The major part of the theory is the behaviour of the exciton resembles that of an atom as its surrounding space shortens. A rather good approximation of an exciton's behaviour is the 3-D model of a particle in a box. The solution of this problem provides a sole mathematical connection between energy states and the dimension of space. Decreasing the volume or the dimensions of the available space, increases the energy of the states. Shown in the diagram is the change in electron energy level and bandgap between nanomaterial and its bulk state.
The following equation shows the relationship between energy level and dimension spacing:
Research results provide an alternative explanation of the shift of properties at nanoscale. In the bulk phase, the surfaces appear to control some of the macroscopically observed properties. However, in nanoparticles, surface molecules do not obey the expected configuration in space. As a result, surface tension changes tremendously.
Classical mechanics view
The Young–Laplace equation can give a background on the investigation of the scale of forces applied to the surface molecules:
Under the assumption of spherical shape and resolving the Young–Laplace equation for the new radii (nm), we estimate the new (GPa). The smaller the radii, the greater the pressure is present. The increase in pressure at the nanoscale results in strong forces toward the interior of the particle. Consequently, the molecular structure of the particle appears to be different from the bulk mode, especially at the surface. These abnormalities at the surface are responsible for changes of inter-atomic interactions and bandgap.
| Physical sciences | Classical mechanics | Physics |
612389 | https://en.wikipedia.org/wiki/Akula-class%20submarine | Akula-class submarine | The Akula class, Soviet designation Project 971 Shchuka-B (, NATO reporting name Akula) is a series of fourth generation nuclear-powered attack submarines (SSNs) first deployed by the Soviet Navy in 1986. There are four sub-classes or flights of Shchuka-B, consisting of the original seven Project 971 boats (codenamed Akula I), commissioned between 1984 and 1990; six Project 971Is (Improved Akulas), commissioned between 1991 and 2009; one Project 971U (Akula II), commissioned in 1995; and one Project 971M (Akula III), commissioned in 2001. The Russians call all of the submarines Shchuka-B, regardless of modifications.
Some confusion may exist as the name Akula (Russian: Акула, meaning "shark" in Russian) was used by the Soviets for a different class of submarines, the Project 941, which is known in the West as the . The Project 971 was named Shchuka-B by the Soviets but given the designation Akula by the West after the name of the lead ship, K-284.
According to defense analyst Norman Polmar, the launch of the first submarine in 1985, "shook everyone [in the West] up", as Western intelligence agencies had not expected the Soviet Union to produce such a boat for another ten years.
Design
The Akula incorporates a double hull system composed of an inner pressure hull and an outer "light" hull. This allows more freedom in the design of the exterior hull shape, resulting in a submarine with more reserve buoyancy than its western analogs.
The distinctive "bulb" or "can" located on top of the Akula's rudder houses its towed sonar array when retracted. Most Akulas have the wake detection system () (SOKS) hydrodynamic sensors, which detect changes in temperature and salinity. They are located on the leading edge of the sail, on the outer hull casing in front of the sail and on the bottom of the hull forward of the sail.
Akulas (excluding Nerpa) are armed with four 533 mm torpedo tubes which can use Type 53 torpedoes or the RPK-2, RPK-6 missile, and four 650 mm torpedo tubes which can use Type 65 torpedoes or the RPK-7 missile. These torpedo tubes are arranged in two rows of four tubes each. The external tubes are mounted outside the pressure hull in one row, above the torpedo tubes, and can only be reloaded in port or with the assistance of a submarine tender. The 650 mm tubes can be fitted with liners to use the 533 mm weaponry. The submarine is also able to use its torpedo tubes to deploy naval mines.
Versions
As with many Soviet/Russian craft, information on the status of the Akula-class submarines is sparse, at best. Information provided by sources varies widely.
Project 971 (Akula I)
Of the seven original Akulas, only three are known to still be in service. These boats are equipped with MGK-540 Skat-3 sonar system (NATO reporting name Shark Gill). The lead boat of the class, K-284 Akula, was decommissioned in 2001, apparently to help save money in the cash-strapped Russian Navy. K-322 Kashalot and K-480 Bars [currently Ak Bars] are in reserve. K-480 Bars was put into reserve in 1998, and was being dismantled in February 2010. Pantera returned to service in January 2008 after a comprehensive overhaul. All were retrofitted with the SOKS hydrodynamic sensors. All submarines before K-391 Bratsk have reactor coolant scoops that are similar to the ones of the SSBNs, long and tubular. Bratsk and subsequent submarines have reactor coolant scoops similar to the short ones on the Oscar IIs (the Typhoon, Akula and Oscar classes use the similar OK-650 reactor).
Project 971 and 971I (Improved Akula I)
The six Akulas of this class are all thought to be in service. They are quieter than the original Akulas. Sources also disagree as to whether construction of this class has been suspended, or if there are a further two units planned.
Improved Akula I Hulls: K-328 Leopard, K-461 Volk, K-154 Tigr, K-419 Kuzbass, K-295 Samara and K-152 Nerpa. These submarines are much quieter than early Akula-class submarines and all have the SOKS hydrodynamic sensors except Leopard.
Project 971U (Akula II)
K-157 Vepr is the only completed Akula II (see the table below). The Akula II is longer and displaces about 700 tons (submerged displacement) more than the Akula I. The added space was used for additional quieting measures. K-157 Vepr became the first Russian submarine that was quieter than the latest U.S. attack submarines of that time, which was the improved (SSN 751 and later). Two of these submarines were used to build the SSBNs.
Project 971M (Akula III)
The K-335 Gepard is the 14th submarine of the class and the only completed Akula III (see the table below) built for the Russian Navy. It was the first submarine commissioned in the Russian Navy since the Kursk disaster, as a result, its commissioning ceremony was an important morale boost for the Russian Navy with President Vladimir Putin in attendance. There is no NATO classification for the Akula III. It is longer and has a larger displacement compared to the Akula II, also it has an enlarged sail and a different towed-array dispenser on the vertical fin. Again, more noise reduction methods were employed. The Gepard was the most advanced Russian submarine before the submarines of the and Borei class were commissioned.
The Soviet advances in sound quieting were of considerable concern to the West, for acoustics was long considered the most significant advantage in U.S. submarine technology compared to the Soviets.
In 1983–1984 the Japanese firm Toshiba sold sophisticated, nine axis milling equipment to the Soviets along with the computer control systems, which were developed by Norwegian firm Kongsberg Vaapenfabrik. U.S Navy officials and Congressmen announced that this technology enabled the Soviet submarine builders to produce more accurate and quieter propellers. This is known as the Toshiba–Kongsberg scandal.
Due to the breakup of the Soviet Union in 1991, production of all Akulas slowed.
The 1999–2000 edition of Jane's Fighting Ships incorrectly listed the first Akula III as Viper (the actual name is "Vepr", "wild boar" in Russian), commissioned on 25 November 1995. Gepard (Cheetah), was launched in 1999 and was commissioned 5 December 2001.
Operational history
Between December 1995 and February 1996, submarine Volk was deployed to the Mediterranean along the Russian aircraft carrier , where she monitored activities of several NATO submarines under Captain 1st rank S. V. Spravtsev.
Between April and June 1996, Tigr was deployed in the Atlantic, where she detected a U.S. SSBN and tracked it on its combat patrol. On 23 July 1996, its commander, Captain 1st rank Alexey Burilichev, received a Hero of the Russian Federation award.
In August 2009, the news media reported that two Akula-class submarines operated off the East Coast of the United States, with one of the submarines being identified as a Project 971 Shchuka-B type. U.S. military sources noted that this was the first known Russian submarine deployment to the western Atlantic since the end of the Cold War, raising concerns within U.S. military and intelligence communities. U.S. Northern Command confirmed that this 2009 Akula-class submarine deployment did occur. One of the boats was likely Gepard that finished a relatively lengthy combat patrol between June and September that year under the command of the Captain 1st rank Alexey Vyacheslavovich Dmitrov, who on 15 February 2012 was awarded a title Hero of the Russian Federation for courage shown at work. The other submarine could have been Tigr under the command of Captain E. A. Petrov, given that she performed a combat patrol sometime between March and November 2009. It is unlikely that other submarines of the project 971 could have been present in the Atlantic that year. Pantera was in Severemorsk during summer, while Vepr, Leopard and Volk did not report any kind of activity in that year (1-3 submarines of the project are usually active with the Northern Fleet at any given moment).
In August 2012, the news media reported that another Akula-class submarine operated in the Gulf of Mexico purportedly undetected for over a month, sparking controversy within U.S. military and political circles, with U.S. Senator John Cornyn of the Senate Armed Services Committee demanding details of this deployment from Admiral Jonathan W. Greenert, the Chief of Naval Operations. Most likely, this was Tigr, as its commander Captain 1st rank Pavel Bulgakov received the Order of Courage on the Defender of the Fatherland Day on 22 February 2013.
Units
Nerpa 2008 accident
On 27 October 2008, it was reported that K-152 Nerpa of the Russian Pacific Fleet had begun her sea trials in the Sea of Japan before handover under a lease agreement to the Indian Navy. On 8 November 2008, while conducting one of these trials, an accidental activation of the halon-based fire-extinguishing system took place in the fore section of the vessel. Within seconds the halon gas had displaced all breathable air from the compartment. As a result, 20 people (17 civilians and 3 seamen) were killed by asphyxiation. Dozens of others suffered freon-related injuries and were evacuated to an unknown port in Primorsky Krai. This was the worst accident in the Russian navy since the loss of the submarine K-141 Kursk in 2000. The submarine itself did not sustain any serious damage and there was no release of radiation.
Lease to India
Chakra II
Three hundred Indian Navy personnel were trained in Russia for the operation of the Akula II submarine Nerpa. India has finalised a deal with Russia, in which at the end of the lease of these submarines, it has an option to buy them. The submarine is named INS Chakra as was the previous India-leased Soviet Charlie-I SSGN. Chakra was officially commissioned into the Indian Navy on 4 April 2012.
Whereas the Russian Navy's Akula-II could be equipped with 28 nuclear-capable cruise missiles with a striking range of , the Indian version is reportedly armed with the -range Club-S nuclear-capable missiles. Missiles with ranges greater than cannot be exported due to arms control restrictions, since Russia is a signatory to the MTCR treaty.
In June 2021, Nerpa was reported in Singapore with Indian crew aboard and on its way back to Russia, despite one year remaining of the 10-year lease, commenced in April 2012. The stated reason was problems with maintenance of the nuclear reactors. Accordingly, the lease will not be prolonged after 2022, as was initially expected.
Chakra III
Russia said in December 2014 that it is ready to lease India more nuclear-powered submarines a day after President Vladimir Putin and Prime Minister Narendra Modi pledged to deepen defence ties.
In January 2015, it was reported that India was involved in negotiations involving the leasing of the Kashalot and the Iribis.
On 7 March 2019, India and Russia signed a $3 billion deal for lease of another Akula-class nuclear-powered attack submarine. The submarine, dubbed as Chakra III, should be delivered to the Indian Navy by 2025.
As of November 2024, the deal was delayed until at least 2028 according to multiple sources.
Gallery
| Technology | Naval warfare | null |
612479 | https://en.wikipedia.org/wiki/Hydrogen%20line | Hydrogen line | The hydrogen line, 21 centimeter line, or H I line is a spectral line that is created by a change in the energy state of solitary, electrically neutral hydrogen atoms. It is produced by a spin-flip transition, which means the direction of the electron's spin is reversed relative to the spin of the proton. This is a quantum state change between the two hyperfine levels of the hydrogen 1 s ground state. The electromagnetic radiation producing this line has a frequency of (1.42 GHz), which is equivalent to a wavelength of in a vacuum. According to the Planck–Einstein relation , the photon emitted by this transition has an energy of []. The constant of proportionality, , is known as the Planck constant.
The hydrogen line frequency lies in the L band, which is located in the lower end of the microwave region of the electromagnetic spectrum. It is frequently observed in radio astronomy because those radio waves can penetrate the large clouds of interstellar cosmic dust that are opaque to visible light. The existence of this line was predicted by Dutch astronomer H. van de Hulst in 1944, then directly observed by E. M. Purcell and his student H. E. Ewen in 1951. Observations of the hydrogen line have been used to reveal the spiral shape of the Milky Way, to calculate the mass and dynamics of individual galaxies, and to test for changes to the fine-structure constant over time. It is of particular importance to cosmology because it can be used to study the early Universe. Due to its fundamental properties, this line is of interest in the search for extraterrestrial intelligence. This line is the theoretical basis of the hydrogen maser.
Cause
An atom of neutral hydrogen consists of an electron bound to a proton. The lowest stationary energy state of the bound electron is called its ground state. Both the electron and the proton have intrinsic magnetic dipole moments ascribed to their spin, whose interaction results in a slight increase in energy when the spins are parallel, and a decrease when antiparallel. The fact that only parallel and antiparallel states are allowed is a result of the quantum mechanical discretization of the total angular momentum of the system. When the spins are parallel, the magnetic dipole moments are antiparallel (because the electron and proton have opposite charge), thus one would expect this configuration to actually have lower energy just as two magnets will align so that the north pole of one is closest to the south pole of the other. This logic fails here because the wave functions of the electron and the proton overlap; that is, the electron is not spatially displaced from the proton, but encompasses it. The magnetic dipole moments are therefore best thought of as tiny current loops. As parallel currents attract, the parallel magnetic dipole moments (i.e., antiparallel spins) have lower energy.
In the ground state, the spin-flip transition between these aligned states has an energy difference of . When applied to the Planck relation, this gives:
where is the wavelength of an emitted photon, is its frequency, is the photon energy, is the Planck constant, and is the speed of light in a vacuum. In a laboratory setting, the hydrogen line parameters have been more precisely measured as:
λ =
ν =
in a vacuum.
This transition is highly forbidden with an extremely small transition rate of , and a mean lifetime of the excited state of around 11 million years. Collisions of neutral hydrogen atoms with electrons or other atoms can help promote the emission of 21 cm photons. A spontaneous occurrence of the transition is unlikely to be seen in a laboratory on Earth, but it can be artificially induced through stimulated emission using a hydrogen maser. It is commonly observed in astronomical settings such as hydrogen clouds in our galaxy and others. Because of the uncertainty principle, its long lifetime gives the spectral line an extremely small natural width, so most broadening is due to Doppler shifts caused by bulk motion or nonzero temperature of the emitting regions.
Discovery
During the 1930s, it was noticed that there was a radio "hiss" that varied on a daily cycle and appeared to be extraterrestrial in origin. After initial suggestions that this was due to the Sun, it was observed that the radio waves seemed to propagate from the centre of the Galaxy. These discoveries were published in 1940 and were noted by Jan Oort who knew that significant advances could be made in astronomy if there were emission lines in the radio part of the spectrum. He referred this to Hendrik van de Hulst who, in 1944, predicted that neutral hydrogen could produce radiation at a frequency of due to two closely spaced energy levels in the ground state of the hydrogen atom.
The 21 cm line (1420.4 MHz) was first detected in 1951 by Ewen and Purcell at Harvard University, and published after their data was corroborated by Dutch astronomers Muller and Oort, and by Christiansen and Hindman in Australia. After 1952 the first maps of the neutral hydrogen in the Galaxy were made, and revealed for the first time the spiral structure of the Milky Way.
Uses
In radio astronomy
The 21 cm spectral line appears within the radio spectrum (in the L band of the UHF band of the microwave window to be exact). Electromagnetic energy in this range can easily pass through the Earth's atmosphere and be observed from the Earth with little interference. The hydrogen line can readily penetrate clouds of interstellar cosmic dust that are opaque to visible light. Assuming that the hydrogen atoms are uniformly distributed throughout the galaxy, each line of sight through the galaxy will reveal a hydrogen line. The only difference between each of these lines is the Doppler shift that each of these lines has. Hence, by assuming circular motion, one can calculate the relative speed of each arm of our galaxy. The rotation curve of our galaxy has been calculated using the hydrogen line. It is then possible to use the plot of the rotation curve and the velocity to determine the distance to a certain point within the galaxy. However, a limitation of this method is that departures from circular motion are observed at various scales.
Hydrogen line observations have been used indirectly to calculate the mass of galaxies, to put limits on any changes over time of the fine-structure constant, and to study the dynamics of individual galaxies. The magnetic field strength of interstellar space can be measured by observing the Zeeman effect on the 21-cm line; a task that was first accomplished by G. L. Verschuur in 1968. In theory, it may be possible to search for antihydrogen atoms by measuring the polarization of the 21-cm line in an external magnetic field.
Deuterium has a similar hyperfine spectral line at 91.6 cm (327 MHz), and the relative strength of the 21 cm line to the 91.6 cm line can be used to measure the deuterium-to-hydrogen (D/H) ratio. One group in 2007 reported D/H ratio in the galactic anticenter to be 21 ± 7 parts per million.
In cosmology
The line is of great interest in Big Bang cosmology because it is the only known way to probe the cosmological "dark ages" from recombination (when stable hydrogen atoms first formed) to the reionization epoch. After including the redshift range for this period, this line will be observed at frequencies from 200 MHz to about 15 MHz on Earth. It potentially has two applications. First, by mapping the intensity of redshifted 21 centimeter radiation it can, in principle, provide a very precise picture of the matter power spectrum in the period after recombination. Second, it can provide a picture of how the universe was re‑ionized, as neutral hydrogen which has been ionized by radiation from stars or quasars will appear as holes in the 21 cm background.
However, 21 cm observations are very difficult to make. Ground-based experiments to observe the faint signal are plagued by interference from television transmitters and the ionosphere, so they must be made from very secluded sites with care taken to eliminate interference. Space based experiments, even on the far side of the Moon (where they would be sheltered from interference from terrestrial radio signals), have been proposed to compensate for this. Little is known about other foreground effects, such as synchrotron emission and free–free emission on the galaxy. Despite these problems, 21 cm observations, along with space-based gravitational wave observations, are generally viewed as the next great frontier in observational cosmology, after the cosmic microwave background polarization.
Relevance to the search for non-human intelligent life
The Pioneer plaque, attached to the Pioneer 10 and Pioneer 11 spacecraft, portrays the hyperfine transition of neutral hydrogen and used the wavelength as a standard scale of measurement. For example, the height of the woman in the image is displayed as eight times 21 cm, or 168 cm. Similarly the frequency of the hydrogen spin-flip transition was used for a unit of time in a map to Earth included on the Pioneer plaques and also the Voyager 1 and Voyager 2 probes. On this map, the position of the Sun is portrayed relative to 14 pulsars whose rotation period circa 1977 is given as a multiple of the frequency of the hydrogen spin-flip transition. It is theorized by the plaque's creators that an advanced civilization would then be able to use the locations of these pulsars to locate the Solar System at the time the spacecraft were launched.
The 21 cm hydrogen line is considered a favorable frequency by the SETI program in their search for signals from potential extraterrestrial civilizations. In 1959, Italian physicist Giuseppe Cocconi and American physicist Philip Morrison published "Searching for interstellar communications", a paper proposing the 21 cm hydrogen line and the potential of microwaves in the search for interstellar communications. According to George Basalla, the paper by Cocconi and Morrison "provided a reasonable theoretical basis" for the then-nascent SETI program. Similarly, Pyotr Makovetsky proposed SETI use a frequency which is equal to either
× ≈
or
2 × ≈
Since is an irrational number, such a frequency could not possibly be produced in a natural way as a harmonic, and would clearly signify its artificial origin. Such a signal would not be overwhelmed by the H I line itself, or by any of its harmonics.
| Physical sciences | Atomic physics | Physics |
613362 | https://en.wikipedia.org/wiki/Solar%20storm | Solar storm | A solar storm is a disturbance on the Sun, which can emanate outward across the heliosphere, affecting the entire Solar System, including Earth and its magnetosphere, and is the cause of space weather in the short-term with long-term patterns comprising space climate.
Types
Solar storms include:
Solar flare, a large explosion in the Sun's atmosphere caused by tangling, crossing or reorganizing of magnetic field lines
Coronal mass ejection (CME), a massive burst of plasma from the Sun, sometimes associated with solar flares
Geomagnetic storm, the interaction of the Sun's outburst with Earth's magnetic field
Solar particle event (SPE), proton or energetic particle (SEP)
| Physical sciences | Solar System | Astronomy |
613480 | https://en.wikipedia.org/wiki/Garter%20snake | Garter snake | Garter snake is the common name for small to medium-sized snakes belonging to the genus Thamnophis in the family Colubridae. They are native to North and Central America, ranging from central Canada in the north to Costa Rica in the south.
With about 35 recognized species and subspecies, garter snakes are highly variable in appearance; generally, they have large round eyes with rounded pupils, a slender build, keeled scales (appearing ‘raised’), and a pattern of longitudinal stripes that may or may not include spots (although some have no stripes at all). Certain subspecies have stripes of blue, yellow, or red, mixed with black tops and beige-tan underbelly markings. They also vary significantly in total length, from .
With no real consensus on the classification of the species of Thamnophis, disagreements between taxonomists and disputed sources (such as field guides) are common. One area of debate, for example, is whether or not two specific types of snake are separate species, or subspecies of the same. Garter snakes are closely related to the genus Nerodia (water snakes), with some species having been moved back and forth between genera.
Garter snakes have been found to be one of the few species of snakes in the world to be both venomous and poisonous. However both are so medically insignificant that it would be extremely unlikely to harm even a human baby. The poisonous aspect comes from their diet so if they are not eating poisonous frogs they are not in fact poisonous as an individual.
Taxonomy
The first garter snake to be scientifically described was the eastern garter snake (now Thamnophis sirtalis sirtalis), by zoologist and taxonomist Carl Linnaeus in 1758. The genus Thamnophis was described by Leopold Fitzinger in 1843 as the genus for the garter snakes and ribbon snakes. Many snakes previously identified as their own genera or species have been reclassified as species or subspecies in Thamnophis. The Reptile Database currently recognised 37 species in the genus, some with several subspecies.
Distribution and habitat
Native to North and Central America, species in the genus Thamnophis can be found in all of the lower 48 United States, and all of the Canadian provinces. They are found from the subarctic plains of west-central Canada east through Ontario and Quebec; from Atlantic Canada and south to Florida, across the southern and central U.S. into the arid regions of the southwest and Mexico, Guatemala and south to the neotropics and Costa Rica.
Garter snakes are not originally native to the eastern Canadian island of Newfoundland, but have been breeding there in the wild and gradually spreading since at least 2010. It is unknown how they reached the island, probably accidentally via hay shipments or as escaped pets.
Their wide distribution is due to their varied diets and adaptability to different habitats, with varying proximity to water. However, in the western part of North America these snakes are more aquatic than in the eastern portion. Garter snakes live in a variety of habitats, including forests, woodlands, fields, grasslands and lawns, but never far from water, often an adjacent wetland, stream or pond. This reflects the fact that amphibians are a large part of their diet. Garter snakes are often found near small ponds with tall weeds.
Behavior
Garter snakes have complex pheromonal communication systems. They can find other snakes by following their pheromone-scented trails. Male and female skin pheromones are so different as to be immediately distinguishable. However, male garter snakes sometimes produce both male and female pheromones. During the mating season, this ability fools other males into attempting to mate with them. This causes the transfer of heat to them in kleptothermy, which is an advantage immediately after hibernation, allowing them to become more active. Male snakes giving off both male and female pheromones have been shown to garner more copulations than normal males in the mating balls that form at the den when females enter the mating melee. A snake hatch can include as many as 57 young.
Garter snakes use the vomeronasal organ to communicate via pheromones through tongue flicking, which gathers chemical cues in the environment. Upon entering the lumen of the organ, the chemical molecules will come into contact with the sensory cells, which are attached to the neurosensory epithelium of the vomeronasal organ.
If disturbed, a garter snake may coil and strike, but it typically hides its head and flails its tail. These snakes will also discharge a malodorous, musky-scented secretion from a gland near the cloaca. This secretion from North American garter snakes contains seven highly odoriferous volatile components: acetic, propanoic, 2-methylpropanoic, butanoic, and 3-methylbutanoic acids; and trimethylamine, and 2-piperidone. They often use these techniques to escape when ensnared by a predator. They will also slither into the water to escape a predator on land. Hawks, crows, egrets, herons, cranes, raccoons, otters and other snake species (such as coral snakes and kingsnakes) will eat garter snakes, with even shrews and frogs eating the juveniles.
Being heterothermic, like all reptiles, garter snakes bask in the sun to regulate their body temperature. During brumation (the reptile equivalent of hibernation), garter snakes typically occupy large communal sites called hibernacula. These snakes will migrate large distances to brumate.
Social behavior
A long-term study by the Ontario Ministry of Transportation has shed light on the social behavior of Butler's garter snakes. The study, conducted in a 250-hectare area near Windsor, Canada, tracked over 3,000 individual snakes over a 12-year period. The findings challenge previous assumptions about solitary snake behavior and suggest that these snakes form social groups and communities. The study revealed that Butler's garter snakes do not wander randomly but instead tend to associate with specific groups of snakes. These groups typically consist of three to four individuals, with some larger groups reaching up to 46 snakes.
Diet
Garter snakes, like all snakes, are carnivorous. Their diet consists of almost any creature they are capable of overpowering: slugs, earthworms (nightcrawlers, as redworms are toxic to garter snakes), leeches, lizards, amphibians (including frog eggs), minnows, and rodents. When living near water, they eat other aquatic animals. The ribbon snake (Thamnophis saurita) in particular favors frogs (including tadpoles), readily eating them despite their strong chemical defenses. Food is swallowed whole. Garter snakes often adapt to eating whatever they can find and whenever they can find it because food can be either scarce or abundant. Although they feed mostly on live animals they will sometimes eat eggs.
Venom
Garter snakes were long thought to be non-venomous, but discoveries in the early 2000s revealed that they produce a neurotoxic venom. Despite this, garter snakes cannot seriously injure or kill humans with the small amounts of comparatively mild venom they produce and they also lack an effective means of delivering it. In a few cases, some swelling and bruising has been reported. They do have enlarged teeth in the back of their mouth but their gums are significantly larger and the secretions of their Duvernoy's gland are only mildly toxic.
Evidence suggests that garter snake and newt populations share an evolutionary link in their tetrodotoxin resistance levels, implying co-evolution between predator and prey. Garter snakes feeding on toxic newts can also retain those toxins in their liver for weeks, making those snakes poisonous as well as venomous.
Conservation status
Despite the decline in their population from collection as pets (especially in the more northerly regions, in which large groups are collected at hibernation), pollution of aquatic areas, and the introduction of American bullfrogs as potential predators, garter snakes are still some of the most commonly found reptiles in much of their ranges. The San Francisco garter snake (Thamnophis sirtalis tetrataenia), however, has been on the endangered list since 1969. Predation by crayfish has also been responsible for the decline of the narrow-headed garter snake (Thamnophis rufipunctatus). Many breeders have bred all species of garter snakes, making it a popular breed.
Species and subspecies
Arranged alphabetically by scientific name:
In the above list, a binomial authority or a trinomial authority in parentheses indicates that the species or subspecies was originally described in a genus other than Thamnophis.
| Biology and health sciences | Reptiles | null |
8841749 | https://en.wikipedia.org/wiki/IPhone | IPhone | The iPhone is a line of smartphones developed and marketed by Apple that run iOS, the company's own mobile operating system. The first-generation iPhone was announced by then–Apple CEO Steve Jobs on January 9, 2007, at Macworld 2007, and launched later that year. Since then, Apple has annually released new iPhone models and iOS versions; the most recent models being the iPhone 16 and 16 Plus, and the higher-end iPhone 16 Pro and 16 Pro Max. As of January 1, 2024, more than 2.3 billion iPhones have been sold, making Apple the largest vendor of mobile phones in 2023.
The original iPhone was the first mobile phone to use multi-touch technology. Throughout its history, the iPhone has gained larger, higher-resolution displays, video-recording functionality, waterproofing, and many accessibility features. Up to the iPhone 8 and 8 Plus, iPhones had a single button on the front panel, with the iPhone 5s and later integrating a Touch ID fingerprint sensor. Since the iPhone X, iPhone models have switched to a nearly bezel-less front screen design with Face ID facial recognition in place of Touch ID for authentication, and increased use of gestures in place of the home button for navigation.
The iPhone, which operates using Apple's proprietary iOS software, is one of the two major smartphone platforms in the world, alongside Android. The first-generation iPhone was described by Steve Jobs as a "revolution" for the mobile phone industry. The iPhone has been credited with popularizing the slate smartphone form factor, and with creating a large market for smartphone apps, or "app economy", laying the foundation for the boom of the market for mobile devices. In addition to the apps that come pre-installed on iOS, there are nearly 2 million apps available for download from Apple's mobile distribution marketplace, the App Store, as of .
History
2000s
Development of an Apple smartphone began in 2004, when the company started to gather a team of 1,000 employees led by hardware engineer Tony Fadell, software engineer Scott Forstall, and design officer Jony Ive, to work on the highly confidential "Project Purple".
Then Apple CEO Steve Jobs steered the original focus away from a tablet (which was later revisited in the form of the iPad) towards a phone. Apple created the device during a secretive collaboration with Cingular Wireless (later renamed AT&T Mobility) at an estimated development cost of US$150 million over thirty months. According to Jobs in 1998, the "i" word in "iMac" (and thereafter "iPod", "iPhone" and "iPad") stands for internet, individual, instruct, inform, and inspire.
Apple rejected the "design by committee" approach that had yielded the Motorola ROKR E1, a largely unsuccessful "iTunes phone" made in collaboration with Motorola. Among other deficiencies, the ROKR E1's firmware limited storage to only 100 iTunes songs to avoid competing with Apple's iPod nano. Cingular gave Apple the liberty to develop the iPhone's hardware and software in-house, a rare practice at the time, and paid Apple a fraction of its monthly service revenue (until the iPhone 3G), in exchange for four years of exclusive U.S. sales, until 2011.
Jobs unveiled the first-generation iPhone to the public on January 9, 2007, at the Macworld 2007 convention at the Moscone Center in San Francisco. The iPhone incorporated a 3.5-inch multi-touch display with few hardware buttons, and ran the iPhone OS operating system with a touch-friendly interface, then marketed as a version of Mac OS X. It was the first mobile phone to use multi-touch technology. The device launched on June 29, 2007, at a starting price of US$499 in the United States, and required a two-year contract with AT&T. The price was reduced by a third after two months. The resulting complaints forced Jobs to issue an apology and offer a partial rebate to early purchasers of the Phone.
On July 11, 2008, at Apple's Worldwide Developers Conference (WWDC) 2008, Apple announced the iPhone 3G, and expanded its launch-day availability to twenty-two countries, and it was eventually released in 70 countries and territories. The iPhone 3G introduced faster 3G connectivity, and a lower starting price of US$199 (with a two-year AT&T contract). It proved commercially popular, overtaking Motorola RAZR V3 as the best selling cell phone in the U.S. by the end of 2008. Its successor, the iPhone 3GS, was announced on June 8, 2009, at WWDC 2009, and introduced video recording functionality.
2010s
The iPhone 4 was announced on June 7, 2010, at WWDC 2010, and introduced a redesigned body incorporating a stainless steel frame and a rear glass panel. At release, the iPhone 4 was marketed as the "world's thinnest smartphone"; it uses the Apple A4 processor, being the first iPhone to use an Apple custom-designed chip. It introduced the Retina display, having four-times the display resolution of preceding iPhones, and was the highest-resolution smartphone screen at release; a front-facing camera was also introduced, enabling video calling functionality via FaceTime.
Users of the iPhone 4 reported dropped/disconnected telephone calls when holding their phones in a certain way, and this issue was nicknamed "antennagate". In January 2011, as Apple's exclusivity agreement with AT&T was expiring, Verizon announced that they would be carrying the iPhone 4, with a model compatible with Verizon's CDMA network releasing on February 10.
The iPhone 4s was announced on October 4, 2011, and introduced the Siri virtual assistant, a dual-core A5 processor, and an 8 megapixel camera with 1080p video recording functionality. The iPhone 5 was announced on September 12, 2012, and introduced a larger 4-inch screen, up from the 3.5-inch screen of all previous iPhone models, as well as faster 4G LTE connectivity. It also introduced a thinner and lighter body made of aluminum alloy, and the 30-pin dock connector of previous iPhones was replaced with the new, reversible Lightning connector.
The iPhone 5s and iPhone 5c were announced on September 10, 2013. The iPhone 5s included a 64-bit A7 processor, becoming the first ever 64-bit smartphone; it also introduced the Touch ID fingerprint authentication sensor. The iPhone 5c was a lower-cost device that incorporated hardware from the iPhone 5, into a series of colorful plastic frames.
On September 9, 2014, Apple introduced the iPhone 6 and iPhone 6 Plus, and included significantly larger screens than the iPhone 5s, at 4.7-inch and 5.5-inch respectively; both models also introduced mobile payment technology via Apple Pay. Optical image stabilization was introduced to the 6 Plus' camera. The Apple Watch was also introduced on the same day, and is a smartwatch that operates in conjunction with a connected iPhone. Some users experienced bending issues from normal use with the iPhone 6 and 6 Plus, particularly on the latter model, and this issue was nicknamed "bendgate".
The iPhone 6s and 6s Plus were introduced on September 9, 2015, and included a more bend-resistant frame made of a stronger aluminum alloy, as well as a higher resolution 12 megapixel main camera capable of 4K video recording. The first-generation iPhone SE was introduced on March 21, 2016, and was a low-cost device that incorporated newer hardware from the iPhone 6s, in the frame of the older iPhone 5s.
The iPhone 7 and 7 Plus were announced on September 7, 2016, which introduced larger camera sensors, IP67-certified water and dust resistance, and a quad-core A10 Fusion processor utilizing big.LITTLE technology; the 3.5 mm headphone jack was removed, and was followed by the introduction of the AirPods wireless earbuds. Optical image stabilization was added to the 7's camera. A second telephoto camera lens was added on the 7 Plus, enabling two-times optical zoom, and "Portrait" photography mode which simulates bokeh in photos.
The iPhone 8, 8 Plus, and iPhone X were announced on September 12, 2017, in Apple's first event held at the Steve Jobs Theater in Apple Park. All models featured rear glass panel designs akin to the iPhone 4, wireless charging, and a hexa-core A11 Bionic chip with "Neural Engine" AI accelerator hardware. The iPhone X additionally introduced a 5.8-inch OLED "Super Retina" display with a "bezel-less" design, with a higher pixel density and contrast ratio than previous iPhones with LCD displays, and introduced a stronger frame made of stainless steel. It also introduced Face ID facial recognition authentication hardware, in a "notch" screen cutout, in place of Touch ID; the home button was removed to achieve the “bezel-less” design, replacing it with a gesture-based navigation system. At its US$999 starting price, the iPhone X was the most expensive iPhone at launch.
The iPhone XR, iPhone XS, and XS Max were announced on September 12, 2018. All models featured the "Smart HDR" computational photography system, and a significantly more powerful "Neural Engine". The XS Max introduced a larger 6.5-inch screen. The iPhone XR included a 6.1-inch LCD "Liquid Retina" display, with a "bezel-less" design similar to the iPhone X, but does not include a second telephoto lens; it was made available in a series of vibrant colors, akin to the iPhone 5c, and was a lower-cost device compared to the iPhone X and XS.
The iPhone 11, 11 Pro, and 11 Pro Max were announced on September 10, 2019. The iPhone 11 was the successor to the iPhone XR, while the iPhone 11 Pro and 11 Pro Max succeeded the iPhone XS and XS Max. All models gained an Ultra-Wide lens, enabling two-times optical zoom out, as well as larger batteries for longer battery life. The second-generation iPhone SE was introduced on April 17, 2020, and was a low-cost device that incorporated newer hardware from the iPhone 11, in the frame of the older iPhone 8, while retaining the home button and the Touch ID sensor.
2020s
The iPhone 12, 12 Mini, 12 Pro, and 12 Pro Max were announced via a livestream event on October 13, 2020. All models featured OLED "Super Retina XDR" displays, introduced faster 5G connectivity, and the MagSafe magnetic charging and accessory system; a slimmer flat-edged design was also introduced, which combined with stronger glass-ceramic front glass, added better drop protection compared to previous iPhones. The iPhone 12 Mini introduced a smaller 5.4-inch screen, while the 12 Pro and 12 Pro Max had larger screens of 6.1-inch and 6.7-inch respectively. The iPhone 12 Pro and 12 Pro Max additionally added a Lidar sensor for better accuracy in augumented reality (AR) applications.
The iPhone 13, 13 Mini, 13 Pro, and 13 Pro Max were announced via a livestream event on September 14, 2021. All models featured larger camera sensors, larger batteries for longer battery life, and a narrower "notch" screen cutout. The iPhone 13 Pro and 13 Pro Max additionally introduced smoother adaptive 120 Hz refresh rate "ProMotion" technology in its OLED display, and three-times optical zoom in the telephoto lens. The low-cost third-generation iPhone SE was introduced on March 8, 2022, and incorporated the A15 Bionic chip from the iPhone 13, but otherwise retained similar hardware to the second-generation iPhone SE.
The iPhone 14, 14 Plus, 14 Pro, and 14 Pro Max were announced on September 7, 2022. All models introduced satellite phone emergency calling functionality. A new 14 Plus model introduced the large 6.7-inch screen size, first seen on the iPhone 12 Pro Max, into a lower-cost device. The iPhone 14 Pro and 14 Pro Max additionally introduced a higher-resolution 48-megapixel main camera, the first increase in megapixel count since the iPhone 6s; it also introduced always-on display technology to the lock screen, and an interactive status bar interface integrated in a redesigned screen cutout, entitled "Dynamic Island".
The iPhone 15, 15 Plus, 15 Pro, and 15 Pro Max were announced on September 12, 2023. Starting with this group of devices, all models switch to using USB-C as their power connector to comply with European Commission regulations, replacing Apple's proprietary Lightning connector after eleven years of use in previous models. The 15 and 15 plus now feature the Dynamic Island, which debuted with the iPhone 14 Pro (effectively retiring the "notch" display cutout), a 48-megapixel main camera, slightly curved edges, and a color-infused frosted glass back. The 15 Pro and Pro Max also replace the mute switch with the "Action" button, and stainless-steel material to titanium.
The iPhone 16, 16 Plus, 16 Pro, and 16 Pro Max were announced on September 9, 2024. The former two introduced a vertical camera layout with refined "Fusion" and Ultra-Wide cameras. The 16 Pro and Pro Max have larger 6.3-inch and 6.9-inch displays, a 48-megapixel Ultra-Wide camera, and the largest batteries in an iPhone up to that point. All models now include access to new Apple Intellegence AI features, a refined thermal system, support for Wi-Fi 7, and a new button dubbed the "Camera Control", allowing easier access to camera features.
Models
46 iPhone models have been produced. The models in bold are devices of the latest generation:
Production
Up to the iPhone 4, all iPhones and other devices, such as iPod Touch models and iPads, were manufactured by Foxconn, based in Taiwan. In 2011, new CEO Tim Cook changed Apple's manufacturing strategy to diversify its supply base. The iPhone 4s in 2012 was the first model to be manufactured simultaneously by two stand-alone companies: Foxconn and Pegatron, the latter also based in Taiwan. Although Foxconn still produces more iPhones, Pegatron's orders have been slowly increased: the company made part of the iPhone 5c line in 2013, and 30% of iPhone 6 devices in 2014. The 6 Plus model was produced solely by Foxconn. In 2019, Apple investigated reports that some Foxconn managers had used rejected parts to build iPhones. In India, Apple pays Wistron, a Taiwan-based manufacturer with a plant near Bangalore, to assemble iPhones to sell in the region.
In 2022, Apple announced that a portion of the iPhone 14 would be manufactured in Tamil Nadu, India, as a response to China's "zero-COVID" policy that has negatively affected global supply chains for many industries. Apple has stated that they plan to shift 25% of iPhone production to India by 2025.
Hardware
Apple directly sub-contracts hardware production to external OEM companies, maintaining a high degree of control over the end product. The iPhone contains most of the hardware parts of a typical modern smartphone. Some hardware elements, such as 3D Touch and the Taptic Engine, are unique to the iPhone. The main hardware of the iPhone is the touchscreen, with current models offering screens of 4.7 inches and larger. All iPhones include a rear-facing camera; the front-facing camera dates back to the iPhone 4. The iPhone 7 Plus introduced multiple lenses to the rear-facing camera. A range of sensors are also included on the device, such as a proximity sensor, ambient light sensor, accelerometer, gyroscopic sensor, magnetometer, facial recognition sensor or fingerprint sensor (depending on the model) and barometer. In 2022, Apple added satellite communications to the iPhone, with the release of the iPhone 14 and iPhone 14 Pro.
Software
Operating system
The iPhone runs iOS. It is based on macOS's Darwin and many of its userland APIs, with Cocoa replaced by Cocoa Touch, and AppKit replaced by UIKit. The graphics stack runs on Metal, Apple's low-level graphics API. The iPhone comes with a set of bundled applications developed by Apple, and supports downloading third-party applications through the App Store.
Apple provides free updates to iOS over-the-air, or through Finder and iTunes on a computer. Major iOS releases have historically accompanied new iPhone models. The most recent version is iOS 18.
App Store and third-party apps
At WWDC 2007 on June 11, 2007, Apple announced that the iPhone would support third-party Ajax web applications that share the look and feel of the iPhone interface. On October 17, 2007, Steve Jobs, in an open letter posted to Apple's "Hot News" weblog, announced that a software development kit (SDK) would be made available to third-party developers in February 2008. The iPhone SDK was officially announced and released on March 6, 2008. The App Store was launched with the release of iPhone OS 2.0, on July 11, 2008.
Apple requires all third-party apps to be downloaded from the App Store, with exceptions for ad-hoc apps used within enterprises. Developers must pay a yearly $99 fee as part of Apple's Developer Program; if their membership expires, their apps are removed from the App Store, though existing users retain the ability to redownload the app. Developers can release free apps, or paid apps for which Apple takes a 30% cut of proceeds. Developers earning less than $1 million in annual sales qualify for the App Store Small Business Program, with Apple only taking a 15% fee.
Though iOS has far lower market share than Android, its app ecosystem has been described as superior, with higher-quality apps, and more iOS-exclusive releases. Android's version fragmentation, less uniform hardware, and lower app revenues have been cited as key factors.
All apps must pass Apple's app review process before being distributed in the App Store. Apple may also stop distributing apps it deems inappropriate. For example, in 2009, Apple rejected the Newspapers app due to The Sun's "obscene" topless Page 3 girls. In 2018, Apple removed Tumblr from the App Store, citing illegal content, causing Tumblr to ban all adult content from their platform. The App Store's review process has been criticized by developers as "frustrating", "anti-competitive", and "asinine".
Users can also install native apps outside of the App Store through jailbreaking, or through exploits, such as TrollStore. Jailbreaking may cause security issues, and is not supported by Apple.
, Apple has passed 60 billion app downloads. , there have been over 140 billion app downloads from the App Store. In January 2017, the App Store had over 2.2million apps for the iPhone. As of August 2024, Apple's App Store contains nearly 2 million applications.
Jailbreaking
Apple restricts the installation of unapproved third-party apps and does not allow full access to the iPhone's filesystem. According to Jonathan Zittrain, the emergence of closed devices like the iPhone has made computing more proprietary than it was in the PC era. Jailbreaking allows users to install apps not available on the App Store, customize their device in ways not allowed by Apple, and bypass SIM locks without carrier approval. Some jailbreak tweaks were later copied by Apple and implemented into iOS, like multitasking, widgets, and copy and paste.
Apple attempted to use the DMCA to fight jailbreaking; however in 2010, the U.S. found jailbreaking to be legal. Jailbroken iPhones are at higher risk of malware due to Apple's lesser control of the app ecosystem. In the United States, Apple cannot void an iPhone's warranty solely due to jailbreaking. Jailbreaks rely on exploits. Apple has improved the iPhone's hardware and software security, making these exploits harder to find; as a result, recent iPhones cannot currently be jailbroken.
Accessibility
The iPhone contains a range of accessibility features to support users' visual, auditory, and motor needs. iPhones can notify users through onscreen banners, audio alerts, vibrations, or the LED flash; vibration patterns can be customized by users. Since iOS 15, Siri can read notifications out loud through earphones, and, since iOS 16, through the device's speakers.
Users with motor needs can use Assistive Touch to customize the way they navigate through menus; it can assist users who have difficulties with some gestures, like pinching, and makes these gestures available by tapping on a menu. The user can create their own gestures and customize the layout of the AssistiveTouch menu. If the user has trouble pressing the Home button, it can be set so that it can be activated with an onscreen tap. Gestures, like rotate and shake, are available even when if the iOS device is mounted on a wheelchair. Head Tracking can be used to control an iPhone using facial movements recognized by the front camera.
Low-vision users can enable VoiceOver, a screen reader which describes what is on the screen, while Siri allows for hands-free interaction. The iPhone also supports wireless braille displays to help users read its interface. Text can be enlarged system-wide. The Magnifier app uses the iPhone's Lidar scanner to identify objects, for example doors, people, and objects, and can describe them to the user, as well as their distance. Door Detection can alert the user through sound, speech, and haptics.
Hearing aids that are part of the Made for iPhone program can be controlled from an iPhone. These hearing aids also feature Live Listen, which enables the iPhone to act as a directional microphone, beaming its audio to compatible hearing aids. Live Listen can help the user hear a conversation in a noisy room or hear someone speaking across the room. Apple built Live Listen support into all AirPods, which can also relay audio from a connected iPhone's microphone. Closed captioning and external TTY devices are supported, while Live Caption can transcribe audio across all apps and display it onscreen. Sound Recognition can recognize surrounding noises, including door bells, kettles, water running, and babies crying, and notify the user with an onscreen alert.
Guided Access helps people with autism, ADHD, or sensory challenges stay focused on a single app. With Guided Access, a parent, teacher, or therapist can limit an iOS device to stay on one app by disabling the Home button and limit the amount of time spent in an app. The user can restrict access to the keyboard or touch input on certain areas of the screen.
Marketing
The original iPhone was heavily promoted before its official announcement, creating buzz and anticipation. Upon its release, it was marketed heavily in television, web and print ads created in partnership with TBWA\Chiat\Day.
Apple's premium market positioning has led the iPhone to be seen as a status symbol.
The Apple ecosystem has been described as a key moat that increases iPhone brand loyalty. iMessage has especially been singled out with its "green bubbles" phenomena. In iMessage, SMS messages from Android users appear as green bubble, rather than the blue bubbles used for texts from other iPhone users. Group chats between iOS and Android are poorly supported; reactions display as text, rather than bubbles, and images are sent through MMS, which degrades image quality. Some teens have described being "ostracized" after switching to Android, which Google has labeled "bullying". This has been described by critics as a key factor leading 87% of U.S. teenagers to use iPhones.
Retail
SIM unlocking
Many iPhones bought through a monthly carrier contract are SIM locked, restricting their use to one particular carrier. While the iPhone was initially sold in the U.S. only on the AT&T network with a SIM lock in place, various hackers found methods to bypass that SIM lock. More than a quarter of first-generation iPhones sold in the U.S. were not registered with AT&T. Apple speculated that they were likely shipped overseas and unlocked, a lucrative market before the iPhone 3G's worldwide release. Today, many carriers either remove the SIM lock automatically after a certain period, or do it upon request, either for free or for a small fee. iPhones bought from Apple are not SIM locked. Many carriers also sell the iPhone unlocked when purchased outright rather than on a long-term contract.
Retail strategy
Since 2013, iPhone buyers can obtain a trade in discount when buying a new iPhone directly from Apple. The program aims to increase the number of customers who purchase iPhones at Apple Stores rather than carrier stores. In 2015, Apple unveiled the iPhone Upgrade Program, a 24-month leasing agreement, which Fortune described as a "change [in] iPhone owners' relationships with mobile carriers".
Repairability
Only Apple Stores and Apple Authorized Service Providers are allowed by Apple to perform genuine replacements. Apple has taken steps to make third-party repairs more difficult. iPhone components are soldered, and many are glued together. iPhones receive low repairability scores, in part due to the difficulty of obtaining genuine parts, and the difficulty undertaking each repair. This has given rise to the right to repair movement, aimed at giving users cheaper options for repairing their phones. Apple has lobbied against right to repair legislation. Multiple jurisdictions aim to introduce right to repair laws, including the EU, UK, and U.S.
In the past, Apple bricked iPhone 6 models after their home buttons were replaced, displaying an Error 53 message; Apple called this a bug, and released an update to address the issue. On iPhones with a Touch ID sensor, the home button cannot be replaced by users or independent repair shops without losing Touch ID functionality, since Apple has not made their calibration tool public.
Starting with the iPhone XR, Apple displays warnings in the Settings app if the battery, display, or camera are replaced by a third party. Additionally, some features are disabled when a part labeled "non-genuine" is detected, like True Tone, or the battery health measurement. iFixit notes that a proprietary, cloud-linked System Configuration tool is required to "complete" a part repair, meaning that even replacing a genuine part with another genuine part will fail Apple's "genuine parts" check unless said tool is used.
In 2022, Apple rolled out a self-service repair program, allowing any user to buy parts, rent repair tools from Apple, and obtain repair manuals. The program received a degree of praise by iFixit and repair advocates, who also critically noted that Apple maintains control over the parts supply.
Privacy
Tracking prevention
Apple introduced App Tracking Transparency (ATT) with iOS 14.5 in April 2021. ATT requires apps to ask for explicit permission before being allowed to track the user across other apps and websites. If the user refuses, the app cannot access Apple's Identifier for Advertisers (IDFA), an identifier used to serve personalized ads. ATT does not prevent personalized ads that are based on the user's behavior within the app itself. The feature has been criticized by some as anti-competitive, including Facebook, whose shares fell by 26% after its rollout. Apple exempts their own apps from their anti-tracking measures, which has led to anti-trust investigations by the French and German governments.
Location tracking controversy
In July 2010, Apple claimed that it collected iPhone users' GPS coordinates and nearby Wi-Fi networks twice a day; a Wall Street Journal investigation found that Google's Android sent this data "several times an hour".
In September 2010, forensic expert Christopher Vance discovered a hidden unencrypted file named "consolidated.db" that contained a record of iPhone users' locations. The file was added with the June 2010 iOS 4 update, though previous versions of iOS stored similar information in a file called "h-cells.plist". On April 20, 2011, The Guardian publicized research by Alasdair Allan and Pete Warden, who found that anyone with physical access to an iPhone could obtain a detailed record of its owner's location and movements over the past year. Moreover, the file was automatically backed up by iTunes onto any computer the iPhone was synchronized with. A Wall Street Journal investigation found that users' locations were still stored when location services are disabled. The controversy led to U.S. congressional scrutiny and an FCC investigation, and was dubbed "Locationgate" by the media.
Apple responded on April 27, 2011, claiming that the data was used to cache nearby Wi-Fi hotspots and cell towers in order to improve location speed and accuracy. The company also claimed that locations being collected when location services were off, and being stored for more than a year, were both bugs. Apple issued an update for iOS (version 4.3.3, or 4.2.8 for the CDMA iPhone 4) which reduced the size of the cache, encrypted it, stopped it being backed up to iTunes, and erased it entirely whenever location services were turned off. Nevertheless, in July 2014, a report on state-owned China Central Television called iPhone tracking a "national security concern".
Currently, iPhones contain a "Frequent Locations" database which records where users have been, along with exact times they arrived and left, raising concerns that the data could be used in court. This feature can be turned off.
Child safety controversy
In August 2021, Apple announced plans to scan iCloud Photos for child abuse imagery (through an algorithm called "NeuralHash"), and filter explicit images sent and received by children using iPhones (dubbed "Conversation Safety"), to be rolled out later that year. More than 90 policy and human rights groups wrote an open letter to condemn both features. Apple's plan to implement NeuralHash on-device rather than in the cloud led the EFF and security experts to call it a "backdoor" that could later be expanded to detect other types of contents, and would decrease users' privacy. Apple claimed the system was "misunderstood", but announced in December 2022 that the photo-scanning feature would never be implemented. The other feature, Conversation Safety, was added in iOS 15.2.
Security
Apple's iOS operating system is regarded by some security experts as more secure against common malware than Android. Less than 1% of mobile malware targets iOS.
Prior to 2014, the iPhone stored all "messages, pictures and videos, contacts, audio recordings [...] and call history" in unencrypted form, enabling easy access by law enforcement. This changed with iOS 8, which adopted file-based encryption. Apple does not hold the decryption key, and cannot be compelled to turn over user data, even when presented with a government warrant. Companies like Grayshift and Cellebrite developed exploits that enable law enforcement to extract user data from iPhones without needing the user's passcode.
In 2015 and 2016, a dispute unfolded between Apple and the FBI. The FBI had recovered the iPhone 5c of one of the San Bernardino attackers, and iCloud backups of that phone from a month and a half before the shooting. The U.S. government attempted to obtain a court order under the All Writs Act compelling Apple to produce a modified version of iOS that would allow investigators to brute force the device passcode. Tim Cook responded on the company's website, outlining a need for encryption, arguing that a backdoor would compromise the privacy of all iPhone users. The DOJ withdrew its request after the FBI bought an exploit to bypass the iPhone's passcode. As a countermeasure, Apple implemented USB Restricted Mode, which was subsequently exploited too.
In 2016,researchers discovered the Pegasus suite of exploits targeting iOS and Android, which led to significant international media coverage. Some Pegasus exploits are zero-click, meaning that they can fully compromise the device with no user interaction, for example by sending a malformed iMessage to the user that would not even trigger a notification. Pegasus can collect most data, including chats, passwords, and photos, and can turn on the phone's microphone and camera remotely. Apple quickly issued an update fixing FORCEDENTRY and other known Pegasus exploits, though Pegasus continued to be used, relying on new exploits. Apple announced a new bug bounty for vulnerabilities, and added an optional Lockdown Mode to iOS 16 that reduces the iPhone's attack surface. Many security researchers have criticized Apple's bug bounty for underpaying researchers, being uncommunicative, and being slow to fix vulnerabilities, and two Apple employees told The Washington Post that the company "has a massive backlog of bugs that it hasn't fixed".
Prominent victims of Pegasus include Jamal Khashoggi, and numerous activists, businessmen and politicians. Pegasus has been widely used since 2011, and is still used by law enforcement and governments as of July 2022.
Reception and legacy
The original iPhone has been described as "revolutionary", a "breakthrough handheld computer", and "the best phone that anybody has ever made". It is now Apple's bestselling product, and has been credited with helping to make Apple one of the world's most valuable publicly traded companies by 2011. Newer iterations have also received praise and awards.
Before the iPhone, smartphones were mostly used for texting, calls, and email; more advanced functions were harder to use and inconvenient on a small screen. They were also hard to develop for, and lacked a thriving app ecosystem like the App Store (released in 2008). Many phones were heavily customized by mobile carriers, which led to feature fragmentation and prevented these phones from turning into thriving software platforms. In contrast, Apple's iPhone SDK provided a wide range of APIs, made mobile development far more accessible, and was instrumental in turning the iPhone into a "Swiss army knife" with a wide range of features and apps.
Successive iPhone models have generated significant fan enthusiasm, with many customers queuing up in front of Apple Stores on launch day. As of 2021, the iPhone has higher brand loyalty than any other smartphone.
The iPhone's success has led to the decline of incumbents Nokia, BlackBerry, and Motorola. RIM, Symbian and Microsoft all attempted to develop more modern operating systems to compete with the iPhone, like Maemo, Windows Phone, and BlackBerry 10; all were unsuccessful. Google successfully started over on their Android project, and designed it for mass adoption by carriers and phone hardware manufacturers. Today, iOS and Android account for 99% of smartphones used worldwide.
Sales
Steve Jobs's initial target was to reach 1% of phone market share in 2008. Apple sold 6.1 million units of the original iPhone between Q3 FY2007 and Q4 FY2008, and 11.3 million units of the iPhone 3G in Q4 FY2008 and Q1 FY2009. In 2008, the iPhone reached 1.1% of worldwide mobile phone market share, and 8.2% of the smartphone market. During this time it was quickly becoming relevant in North America, and in market share was ranked second in the U.S. in 2009, behind the BlackBerry; in 2010 the iPhone 3GS was the best-selling smartphone in the U.S., the first time that an iPhone device reached top spot in that market.
iPhone sales grew continuously year-over-year since its introduction until Q2 FY2016. The iPhone briefly surpassed BlackBerry in Q4 FY2008, and permanently overtook it starting in Q3 FY2010. By 2011, Apple sold 100 million iPhones worldwide, and became the largest mobile phone vendor in the world by revenue, surpassing long-time leader Nokia. Q1 FY2012 marked Apple's best quarterly earnings in its history, with 53% of the company's revenues coming from iPhone sales. Phone sales are strongly seasonal, peaking in the holiday season (Apple's Q1). With the release of the iPhone 13 in Q1 FY2022, Apple temporarily topped Samsung, with 84.9 million units shipped compared to Samsung's 68.9 million. In most quarters, Apple is the second largest smartphone vendor by units. Apple sold 223 million iPhones in its financial year 2023 ending September 24.
Today, Samsung and Apple dominate the smartphone market, with 21.8% and 15.6% worldwide market share respectively. Due to Apple's small lineup, Apple often dominates the list of bestselling smartphone models. Despite its lower market share, the iPhone's premium positioning has led it to capture nearly half of global smartphone revenue, and 80% of global smartphone profits, with Samsung taking the other 20%. Carriers compete with each other to subsidize iPhone upgrades, which is seen as a significant factor in iPhone sales, though this has reduced carrier profits. On July 27, 2016, Apple announced that it had sold their 1 billionth iPhone. As of January 1, 2024, more than 2.3 billion iPhones have been sold.
Compared to other high-tech products, a greater proportion of iPhone users are female. The iPhone has been adopted by both consumers and business users. iPhone users are wealthier and spend more time on their phones than Android users on average. The iPhone is especially popular in the U.S., where it has a 50% market share, and is used by 87% of teenagers. Worldwide, the iPhone accounts for 78% of the high-end ($1,000+) smartphone market.
Android overtook the iPhone's installed base in 2010, according to NPD Group. During Apple's earnings call on January 27, 2021, Tim Cook said that 1 billion iPhones were being actively used worldwide.
Emerging markets
While other manufacturers make separate entry-level phones, Apple's entry-level phones are the previous years' models, part of an effort to increase its market share in emerging markets without diluting its premium brand. It also considers emerging market tastes in its product designs; for example, it introduced a gold iPhone after finding that gold was seen as a popular sign of a luxury product among Chinese customers. In 2017, Apple started manufacturing previous years' iPhone models in India; in 2022, it began manufacturing the current iPhone 14 there too. Analysts have speculated that this was partly caused by Apple's desire to reduce its dependence on China, and to overcome Indian import duties. In 2023, the Chinese government banned the use of iPhones by government civil servants in what was seen as an effort to reduce dependence on foreign technology and strengthen cybersecurity.
In May 2024 Iranian president Mokhber banned imported iPhone 14 and newer models, in November the ban was lifted and replaced with 30% customs tariff to the phones.
| Technology | Computer hardware | null |
1590842 | https://en.wikipedia.org/wiki/Work%20hardening | Work hardening | Work hardening, also known as strain hardening, is the process by which a material's load-bearing capacity (strength) increases during plastic (permanent) deformation. This characteristic is what sets ductile materials apart from brittle materials. Work hardening may be desirable, undesirable, or inconsequential, depending on the application.
This strengthening occurs because of dislocation movements and dislocation generation within the crystal structure of the material. Many non-brittle metals with a reasonably high melting point as well as several polymers can be strengthened in this fashion. Alloys not amenable to heat treatment, including low-carbon steel, are often work-hardened. Some materials cannot be work-hardened at low temperatures, such as indium, however others can be strengthened only via work hardening, such as pure copper and aluminum.
Undesirable work hardening
An example of undesirable work hardening is during machining when early passes of a cutter inadvertently work-harden the workpiece surface, causing damage to the cutter during the later passes. Certain alloys are more prone to this than others; superalloys such as Inconel require materials science machining strategies that take it into account.
For metal objects designed to flex, such as springs, specialized alloys are usually employed in order to avoid work hardening (a result of plastic deformation) and metal fatigue, with specific heat treatments required to obtain the necessary characteristics.
Intentional work hardening
An example of desirable work hardening is that which occurs in metalworking processes that intentionally induce plastic deformation to exact a shape change. These processes are known as cold working or cold forming processes. They are characterized by shaping the workpiece at a temperature below its recrystallization temperature, usually at ambient temperature. Cold forming techniques are usually classified into four major groups: squeezing, bending, drawing, and shearing. Applications include the heading of bolts and cap screws and the finishing of cold rolled steel. In cold forming, metal is formed at high speed and high pressure using tool steel or carbide dies. The cold working of the metal increases the hardness, yield strength, and tensile strength.
Theory
Before work hardening, the lattice of the material exhibits a regular, nearly defect-free pattern (almost no dislocations). The defect-free lattice can be created or restored at any time by annealing. As the material is work hardened it becomes increasingly saturated with new dislocations, and more dislocations are prevented from nucleating (a resistance to dislocation-formation develops). This resistance to dislocation-formation manifests itself as a resistance to plastic deformation; hence, the observed strengthening.
In metallic crystals, this is a reversible process and is usually carried out on a microscopic scale by defects called dislocations, which are created by fluctuations in local stress fields within the material culminating in a lattice rearrangement as the dislocations propagate through the lattice. At normal temperatures the dislocations are not annihilated by annealing. Instead, the dislocations accumulate, interact with one another, and serve as pinning points or obstacles that significantly impede their motion. This leads to an increase in the yield strength of the material and a subsequent decrease in ductility.
Such deformation increases the concentration of dislocations which may subsequently form low-angle grain boundaries surrounding sub-grains. Cold working generally results in a higher yield strength as a result of the increased number of dislocations and the Hall–Petch effect of the sub-grains, and a decrease in ductility. The effects of cold working may be reversed by annealing the material at high temperatures where recovery and recrystallization reduce the dislocation density.
A material's work hardenability can be predicted by analyzing a stress–strain curve, or studied in context by performing hardness tests before and after a process.
Elastic and plastic deformation
Work hardening is a consequence of plastic deformation, a permanent change in shape. This is distinct from elastic deformation, which is reversible. Most materials do not exhibit only one or the other, but rather a combination of the two. The following discussion mostly applies to metals, especially steels, which are well studied. Work hardening occurs most notably for ductile materials such as metals. Ductility is the ability of a material to undergo plastic deformations before fracture (for example, bending a steel rod until it finally breaks).
The tensile test is widely used to study deformation mechanisms. This is because under compression, most materials will experience trivial (lattice mismatch) and non-trivial (buckling) events before plastic deformation or fracture occur. Hence the intermediate processes that occur to the material under uniaxial compression before the incidence of plastic deformation make the compressive test fraught with difficulties.
A material generally deforms elastically under the influence of small forces; the material returns quickly to its original shape when the deforming force is removed. This phenomenon is called elastic deformation. This behavior in materials is described by Hooke's Law. Materials behave elastically until the deforming force increases beyond the elastic limit, which is also known as the yield stress. At that point, the material is permanently deformed and fails to return to its original shape when the force is removed. This phenomenon is called plastic deformation. For example, if one stretches a coil spring up to a certain point, it will return to its original shape, but once it is stretched beyond the elastic limit, it will remain deformed and won't return to its original state.
Elastic deformation stretches the bonds between atoms away from their equilibrium radius of separation, without applying enough energy to break the inter-atomic bonds. Plastic deformation, on the other hand, breaks inter-atomic bonds, and therefore involves the rearrangement of atoms in a solid material.
Dislocations and lattice strain fields
In materials science parlance, dislocations are defined as line defects in a material's crystal structure. The bonds surrounding the dislocation are already elastically strained by the defect compared to the bonds between the constituents of the regular crystal lattice. Therefore, these bonds break at relatively lower stresses, leading to plastic deformation.
The strained bonds around a dislocation are characterized by lattice strain fields. For example, there are compressively strained bonds directly next to an edge dislocation and strained in tension bonds beyond the end of an edge dislocation. These form compressive strain fields and tensile strain fields, respectively. Strain fields are analogous to electric fields in certain ways. Specifically, the strain fields of dislocations obey similar laws of attraction and repulsion; in order to reduce overall strain, compressive strains are attracted to tensile strains, and vice versa.
The visible (macroscopic) results of plastic deformation are the result of microscopic dislocation motion. For example, the stretching of a steel rod in a tensile tester is accommodated through dislocation motion on the atomic scale.
Increase of dislocations and work hardening
Increase in the number of dislocations is a quantification of work hardening. Plastic deformation occurs as a consequence of work being done on a material; energy is added to the material. In addition, the energy is almost always applied fast enough and in large enough magnitude to not only move existing dislocations, but also to produce a great number of new dislocations by jarring or working the material sufficiently enough. New dislocations are generated in proximity to a Frank–Read source.
Yield strength is increased in a cold-worked material. Using lattice strain fields, it can be shown that an environment filled with dislocations will hinder the movement of any one dislocation. Because dislocation motion is hindered, plastic deformation cannot occur at normal stresses. Upon application of stresses just beyond the yield strength of the non-cold-worked material, a cold-worked material will continue to deform using the only mechanism available: elastic deformation, the regular scheme of stretching or compressing of electrical bonds (without dislocation motion) continues to occur, and the modulus of elasticity is unchanged. Eventually the stress is great enough to overcome the strain-field interactions and plastic deformation resumes.
However, ductility of a work-hardened material is decreased. Ductility is the extent to which a material can undergo plastic deformation, that is, it is how far a material can be plastically deformed before fracture. A cold-worked material is, in effect, a normal (brittle) material that has already been extended through part of its allowed plastic deformation. If dislocation motion and plastic deformation have been hindered enough by dislocation accumulation, and stretching of electronic bonds and elastic deformation have reached their limit, a third mode of deformation occurs: fracture.
Quantification of work hardening
The shear strength, , of a dislocation is dependent on the shear modulus, G, the magnitude of the Burgers vector, b, and the dislocation density, :
where is the intrinsic strength of the material with low dislocation density and is a correction factor specific to the material.
As shown in Figure 1 and the equation above, work hardening has a half root dependency on the number of dislocations. The material exhibits high strength if there are either high levels of dislocations (greater than 1014 dislocations per m2) or no dislocations. A moderate number of dislocations (between 107 and 109 dislocations per m2) typically results in low strength.
Example
For an extreme example, in a tensile test a bar of steel is strained to just before the length at which it usually fractures. The load is released smoothly and the material relieves some of its strain by decreasing in length. The decrease in length is called the elastic recovery, and the result is a work-hardened steel bar. The fraction of length recovered (length recovered/original length) is equal to the yield-stress divided by the modulus of elasticity. (Here we discuss true stress in order to account for the drastic decrease in diameter in this tensile test.) The length recovered after removing a load from a material just before it breaks is equal to the length recovered after removing a load just before it enters plastic deformation.
The work-hardened steel bar has a large enough number of dislocations that the strain field interaction prevents all plastic deformation. Subsequent deformation requires a stress that varies linearly with the strain observed, the slope of the graph of stress vs. strain is the modulus of elasticity, as usual.
The work-hardened steel bar fractures when the applied stress exceeds the usual fracture stress and the strain exceeds usual fracture strain. This may be considered to be the elastic limit and the yield stress is now equal to the fracture toughness, which is much higher than a non-work-hardened steel yield stress.
The amount of plastic deformation possible is zero, which is less than the amount of plastic deformation possible for a non-work-hardened material. Thus, the ductility of the cold-worked bar is reduced.
Substantial and prolonged cavitation can also produce strain hardening.
Empirical relations
There are two common mathematical descriptions of the work hardening phenomenon. Hollomon's equation is a power law relationship between the stress and the amount of plastic strain:
where σ is the stress, K is the strength index or strength coefficient, εp is the plastic strain and n is the strain hardening exponent. Ludwik's equation is similar but includes the yield stress:
If a material has been subjected to prior deformation (at low temperature) then the yield stress will be increased by a factor depending on the amount of prior plastic strain ε0:
The constant K is structure dependent and is influenced by processing while n is a material property normally lying in the range 0.2–0.5. The strain hardening index can be described by:
This equation can be evaluated from the slope of a log(σ) – log(ε) plot. Rearranging allows a determination of the rate of strain hardening at a given stress and strain:
Work hardening in specific materials
Steel
Steel is an important engineering material, used in many applications. Steel may be work hardened by deformation at low temperature, called cold working. Typically, an increase in cold work results in a decrease in the strain hardening exponent. Similarly, high strength steels tend to exhibit a lower strain hardening exponent.
Copper
Copper was the first metal in common use for tools and containers since it is one of the few metals available in non-oxidized form, not requiring the smelting of an ore. Copper is easily softened by heating and then cooling (it does not harden by quenching, e.g., quenching in cool water). In this annealed state it may then be hammered, stretched and otherwise formed, progressing toward the desired final shape but becoming harder and less ductile as work progresses. If work continues beyond a certain hardness the metal will tend to fracture when worked and so it may be re-annealed periodically as shaping continues. Annealing is stopped when the workpiece is near its final desired shape, and so the final product will have a desired strength and hardness. The technique of repoussé exploits these properties of copper, enabling the construction of durable jewelry articles and sculptures (such as the Statue of Liberty).
Gold and other precious metals
Much gold jewelry is produced by casting, with little or no cold working; which, depending on the alloy grade, may leave the metal relatively soft and bendable. However, a jeweler may intentionally use work hardening to strengthen wearable objects that are exposed to stress, such as rings.
Aluminum
Items made from aluminum and its alloys must be carefully designed to minimize or evenly distribute flexure, which can lead to work hardening and, in turn, stress cracking, possibly causing catastrophic failure. For this reason modern aluminum aircraft will have an imposed working lifetime (dependent upon the type of loads encountered), after which the aircraft must be retired.
| Technology | Metallurgy | null |
1590884 | https://en.wikipedia.org/wiki/Chinese%20garden | Chinese garden | The Chinese garden is a landscape garden style which has evolved over three thousand years. It includes both the vast gardens of the Chinese emperors and members of the imperial family, built for pleasure and to impress, and the more intimate gardens created by scholars, poets, former government officials, soldiers and merchants, made for reflection and escape from the outside world. They create an idealized miniature landscape, which is meant to express the harmony that should exist between man and nature.
The art of Chinese garden integrates architecture, calligraphy and painting, sculpture, literature, gardening and other arts. It is a model of Chinese aesthetics, reflecting the profound philosophical thinking and pursuit of life of the Chinese people. Among them, Chengde Mountain Resort and the Summer Palace, which belong to royal gardens, and several the Classical Gardens of Suzhou in Jiangsu Province, which belong to private gardens, are also included in the World Heritage List by UNESCO. Many essential elements are used in Chinese gardens, and Moon Gate is one of them.
A typical Chinese garden is enclosed by walls and includes one or more ponds, rock works, trees and flowers, and an assortment of halls and pavilions within the garden, connected by winding paths and zig-zag galleries. By moving from structure to structure, visitors can view a series of carefully composed scenes, unrolling like a scroll of landscape paintings.
History
Beginnings
The earliest recorded Chinese gardens were created in the valley of the Yellow River, during the Shang dynasty (1600–1046 BC). These gardens were large enclosed parks where the kings and nobles hunted game, or where fruit and vegetables were grown. Early inscriptions from this period, carved on tortoise shells, have three Chinese characters for garden, you, pu and yuan. You was a royal garden where birds and animals were kept, while pu was a garden for plants. During the Qin dynasty (221–206 BC), yuan became the character for all gardens. The old character for yuan is a small picture of a garden; it is enclosed in a square which can represent a wall, and has symbols which can represent the plan of a structure, a small square which can represent a pond, and a symbol for a plantation or a pomegranate tree.
A famous royal garden of the late Shang dynasty was the Terrace, Pond and Park of the Spirit (Lingtai, Lingzhao Lingyou) built by King Wenwang west of his capital city, Yin. The park was described in the Classic of Poetry this way:
The King makes his promenade in the Park of the Spirit,
The deer are kneeling on the grass, feeding their fawns,
The deer are beautiful and resplendent.
The immaculate cranes have plumes of a brilliant white.
The King makes his promenade to the Pond of the Spirit,
The water is full of fish, who wriggle.
Another early royal garden was Shaqui, or the Dunes of Sand, built by the last Shang ruler, King Zhou (1075–1046 BC). It was composed of an earth terrace, or tai, which served as an observation platform in the center of a large square park. It was described in one of the early classics of Chinese literature, the Records of the Grand Historian (Shiji). According to the Shiji, one of the most famous features of this garden was the Wine Pool and Meat Forest (酒池肉林). A large pool, big enough for several small boats, was constructed on the palace grounds, with inner linings of polished oval shaped stones from the seashore. The pool was then filled with wine. A small island was constructed in the middle of the pool, where trees were planted, which had skewers of roasted meat hanging from their branches. King Zhou and his friends and concubines drifted in their boats, drinking the wine with their hands and eating the roasted meat from the trees. Later Chinese philosophers and historians cited this garden as an example of decadence and bad taste.
During the Spring and Autumn period (722–481 BC), in 535 BC, the Terrace of Shanghua, with lavishly decorated palaces, was built by King Jing of the Zhou dynasty. In 505 BC, an even more elaborate garden, the Terrace of Gusu, was begun. It was located on the side of a mountain, and included a series of terraces connected by galleries, along with a lake where boats in the form of blue dragons navigated. From the highest terrace, a view extended as far as Lake Tai, the Great Lake.
The Legend of the Isle of the Immortals
An ancient Chinese legend played an important part in early garden design. In the 4th century BC, a tale in the Classic of Mountains and Seas described a peak called Mount Penglai located on one of three islands at the eastern end of the Bohai Sea, between China and Korea, which was the home of the Eight Immortals. On this island were palaces of gold and silver, with jewels on the trees. There was no pain, no winter, wine glasses and rice bowls were always full, and fruits, when eaten, granted eternal life.
In 221 BC, Ying Zheng, the King of Qin conquered other rival states and unified China under the Qin Empire, which he ruled until 210 BC. He heard the legend of the islands and sent emissaries to find the islands and bring back the elixir of immortal life, without success. At his palace near his capital, Xianyang, he created a garden with a large lake called Lanchi gong or the Lake of the Orchids. On an island in the lake he created a replica of Mount Penglai, symbolizing his search for paradise. After his death, the Qin Empire fell in 206 BC and his capital city and garden were completely destroyed, but the legend continued to inspire Chinese gardens. Some gardens have a single island with an artificial mountain representing the island of the Eight Immortals. Other gardens have gardens featuring three Boshan Mountains - Penglai, Yingzhou, and Fanghu or Fangzhang. The Yichi Sanshan () system of one pond with three mountains has been a main model of royal gardens.
Han dynasty (206 BC–220 AD)
Under the Han dynasty (206 BC – 220 AD), a new imperial capital was built at Chang'an, and Emperor Wu built a new imperial garden, which combined the features of botanical and zoological gardens, as well as the traditional hunting grounds. Inspired by another version of Chinese classic about the Isles of the Immortals, called Liezi, he created a large artificial lake, the Lake of the Supreme Essence, with three artificial islands in the center representing the three isles of the Immortals - Penglai,Fanghu, and Yingzhou. The park was later destroyed, but its memory would continue to inspire Chinese garden design for centuries. The Jianzhang Palace in the Han Dynasty is the first known garden built with the complete set of the three remaining Bohai Shenshan mountains. Since then, the Yichi Sanshan () system of one pond with three mountains has been a main model of royal gardens.
Another notable garden of the Han period was the Garden of General Liang Ji built under Emperor Shun (125–144 AD). Using a fortune amassed during his twenty years in the imperial court, Liang Ji built an immense landscape garden with artificial mountains, ravines and forests, filled with rare birds and domesticated wild animals. This was one of the first gardens that tried to create an idealized copy of nature.
Gardens for poets and scholars (221–618 AD)
After the fall of the Han dynasty, a long period of political instability began in China. Buddhism was introduced into China by Emperor Ming (57–75 AD), and spread rapidly. By 495, the city of Luoyang, capital of the Northern Wei dynasty, had over 1,300 temples, mostly in the former residences of believers. Each of the temples had its own small garden.
During this period, many former government officials left the court and built gardens where they could escape the outside world and concentrate on nature and literature. One example was the Jingu Yuan, or Garden of the Golden Valley, built in 296 by Shi Chong (249–300 AD), an aristocrat and former court official, ten kilometers northeast of Luoyang. He invited thirty famous poets to a banquet in his garden, and wrote about the event himself:
This visit to the garden resulted in a famous collection of poems, Jingu Shi, or Poems of the Golden Valley, and launched a long tradition of writing poetry in and about gardens.
The poet and calligrapher Wang Xizhi (307–365) wrote in his excellent calligraphy the Preface to the Poems Composed at the Orchid Pavilion introducing a book recording the event of the Orchid Pavilion Gathering, another famous poetry setting at a country retreat called the "Orchid Pavilion". This was a park with a meandering stream. He brought together a group of famous poets, and seated them beside the stream. Then he placed cups of wine in the stream, and let them float. If the cup stopped beside one of the poets, he was obliged to drink it and then compose a poem. The garden of the floating cup (liubei tang), with small pavilions and artificial winding streams, became extremely popular in both imperial and private gardens.
The Orchid Pavilion inspired Emperor Yang (604–617) of the Sui dynasty to build his new imperial garden, the Garden of the West, near Hangzhou. His garden had a meandering stream for floating glasses of wine and pavilions for writing poetry. He also used the park for theatrical events; he launched small boats on his stream with animated figures illustrating the history of China.
Tang dynasty (618–907), First Golden Age of the Classical Garden
The Tang dynasty (618–907 AD) was considered the first golden age of the classical Chinese garden. Emperor Xuanzong built a magnificent imperial garden, the Garden of the Majestic Clear Lake, near Xi′an, and lived there with his famous concubine, Consort Yang.
Painting and poetry reached a level never seen before, and new gardens, large and small, filled the capital city, Chang'an. The new gardens, were inspired by classical legends and poems. There were shanchi yuan, gardens with artificial mountains and ponds, inspired by the legend of the isles of immortals, and shanting yuan, gardens with replicas of mountains and small viewing houses, or pavilions. Even ordinary residences had tiny gardens in their courtyards, with terracotta mountains and small ponds.
These Chinese classical gardens, or scholar's gardens (wenren yuan), were inspired by, and in turn inspired, classical Chinese poetry and painting. A notable example was the Jante Valley Garden of the poet-painter and civil servant Wang Wei (701–761). He bought the ruined villa of a poet, located near the mouth of a river and a lake. He created twenty small landscape scenes within his garden, with names such as the Garden of Magnolias, the Waving Willows, the Kiosk in the Heart of the Bamboos, the Spring of the Golden Powder, and the View-House beside the Lake. He wrote a poem for each scene in the garden and commissioned a famous artist, to paint scenes of the garden on the walls of his villa. After retiring from the government, he passed his time taking boat trips on the lake, playing the cithare and writing and reciting poetry.
During the Tang dynasty, plant cultivation was developed to an advanced level, with many plant species being grown by means of plant introduction, domestication, transplantation, and grafting. The aesthetic properties of plants were highlighted, while numerous books on plant classification and cultivation were published. The capital, Chang'an, was a very cosmopolitan city, filled with diplomats, merchants, pilgrims, monks and students, who carried descriptions of the gardens all over Asia. The economic prosperity of the Tang dynasty led to the increasing construction of classical gardens across all of China.
The last great garden of the Tang dynasty was the Hamlet of the Mountain of the Serene Spring (Pingquan Shanzhuang), built east of the city of Luoyang by Li Deyu, Grand Minister of the Tang Empire. The garden was vast, with over a hundred pavilions and structures, but it was most famous for its collection of exotically shaped rocks and plants, which its creator collected all over China. Rocks of unusual shapes, known as Chinese Scholars' Rocks, often selected to portray the part of a mountain or mountain range in a garden scene, gradually became an essential feature of the Chinese garden.
Song Dynasty (960–1279)
There were two periods of the Song dynasty, northern and southern, and both were known for the construction of famous gardens. Emperor Huizong (1082–1135) was an accomplished painter of birds and flowers. A scholar himself, he integrated elements of the scholar garden into his grand imperial garden. His first garden, called The Basin of the Clarity of Gold, was an artificial lake surrounded by terraces and pavilions. The public was invited into the garden in the spring for boat races and spectacles on the lake. In 1117 he personally supervised the building of a new garden. He had exotic plants and picturesque rocks brought from around China for his garden, particularly the prized rocks from Lake Tai. Some of the rocks were so large that, in order to move them by water on the grand canal, he had to destroy all the bridges between Hangzhou and Beijing. In the center of his garden he had constructed an artificial mountain a hundred meters high, with cliffs and ravines, which he named Genyue, or "The Mountain of Stability." The garden was finished in 1122. In 1127, Emperor Huizong was forced to flee from the Song capital, Kaifeng, when it came under attack by the armies of the Jurchen-led Jin dynasty. When he returned (as a captive of the Jurchens), he found his garden completely destroyed, all the pavilions burned and the art works looted. Only the mountain remained.
While the imperial gardens were the best known, many smaller but equally picturesque gardens were built in cities such as Luoyang. The Garden of the Monastery of the Celestial Rulers in Luoyang was famous for its peonies; the entire city came when they were in bloom. The Garden of Multiple Springtimes was famous for its view of the mountains. The most famous garden in Luoyang was The Garden of Solitary Joy (Dule Yuan), built by the poet and historian Sima Guang (1021–1086). His garden had an area of eight mu, or about 1.5 hectares. In the center was the Pavilion of Study, his library, with five thousand volumes. To the north was an artificial lake, with a small island, with a picturesque fisherman's hut. To the east was a garden of medicinal herbs, and to the west was an artificial mountain, with a belevedere at the summit to view the surrounding neighborhoods. Any passer-by could visit the garden by paying a small fee.
After fall of Kaifeng, the capital of the Song dynasty was moved to Lin'an (present-day Hangzhou, Zhejiang). The city of Lin'an soon had more than fifty gardens built on the shore of the Western Lake. The other city in the province famous for its gardens was Suzhou, where many scholars, government officials and merchants built residences with gardens. Some of these gardens still exist today, though most have been greatly altered over the centuries.
The oldest Suzhou garden that can be seen today is the Blue Wave Pavilion, built in 1044 by the Song dynasty poet Su Shunqing. (1008–1048). In the Song dynasty, it consisted of a hilltop viewing pavilion. Other lakeside pavilions were added, including a reverence hall, a recitation hall, and a special pavilion for watching the fish. Over the centuries it was much modified, but still keeps its essential plan.
Another Song dynasty garden still in existence is the Master of the Nets Garden in Suzhou. It was created in 1141 by Shi Zhengzhi, Deputy Civil Service Minister of the Southern Song government. It had his library, the Hall of Ten Thousand Volumes, and an adjacent garden called the Fisherman's Retreat. It was extensively remodeled between 1736 and 1796, but it remains one of the best example of a Song Dynasty Scholars Garden.
In the city of Wuxi, on the edge of Lake Tai and at the foot of two mountains, there were thirty four gardens recorded by the Song dynasty historian Zhou Mi (1232–1308). The two most famous gardens, the Garden of the North (Beiyuan) and the Garden of the South (Nanyuan), both belonged to Shen Dehe, Grand Minister to Emperor Gaozong (1131–1162). The Garden of the South was a classic mountain-and-lake (shanshui) garden; it had a lake with an Island of Immortality (Penglai dao), on which were three great boulders from Taihu. The Garden of the South was a water garden, with five large lakes connected to Lake Tai. A terrace gave visitors a view of the lake and the mountains.
Yuan dynasty (1279–1368)
In 1271, Kublai Khan established the Mongol-led Yuan dynasty in China. By 1279, he annihilated the last resistance of the Song dynasty and unified China under Mongol rule. He established a new capital on the site of present-day Beijing, called Dadu, the Great Capital.
The most famous garden of the Yuan dynasty was Kublai Khan's summer palace and garden at Xanadu. The Venetian traveler Marco Polo is believed to have visited Xanadu in about 1275, and described the garden this way:
"Round this Palace a wall is built, inclosing a compass of 16 miles, and inside the Park there are fountains and rivers and brooks, and beautiful meadows, with all kinds of wild animals (excluding such as are of ferocious nature), which the Emperor has procured and placed there to supply food for his gerfalcons and hawks, which he keeps there in mew. Of these there are more than 200 gerfalcons alone, without reckoning the other hawks. The Khan himself goes every week to see his birds sitting in mew, and sometimes he rides through the park with a leopard behind him on his horse's croup; and then if he sees any animal that takes his fancy, he slips his leopard at it, and the game when taken is made over to feed the hawks in mew. This he does for diversion."
This brief description later inspired the poem Kubla Khan by the English romantic poet, Samuel Taylor Coleridge.
When he established his new capital at Dadu, Kublai Khan enlarged the artificial lakes that had been created a century earlier by the Jurchen-led Jin dynasty, and built up the island of Oinghua, creating a striking contrast between curving banks of the lake and garden and the strict geometry of what later became the Forbidden City of Beijing. This contrast is still visible today.
Despite the Mongol invasion, the classical Chinese scholar's garden continued to flourish in other parts of China. An excellent example was the Lion Grove Garden in Suzhou. It was built in 1342, and took its name from the collection of fantastic and grotesque assemblies of rocks, taken from Lake Tai. Some of them were said to look like the heads of lions. The Kangxi and Qianlong emperors of the Qing dynasty each visited the garden several times, and used it as model for their own summer garden, the Garden of Perfect Splendor, at the Chengde Mountain Resort.
In 1368, forces of the Ming dynasty, led by Zhu Yuanzhang, captured Dadu from the Mongols and overthrew the Yuan dynasty. Zhu Yuanzhang ordered the Yuan palaces in Dadu to be burned down.
Ming dynasty (1368–1644)
The most famous existing garden from the Ming dynasty is the Humble Administrator's Garden in Suzhou. It was built during the reign of the Zhengde Emperor (1506–1521) by Wang Xianchen, a minor government administrator who retired from government service and devoted himself to his garden. The garden has been much altered since it was built, but the central part has survived; a large pond full of lotus blossoms, surrounded by structures and pavilions designed as viewpoints of the lake and gardens. The park has an island, the Fragrant Isle, shaped like a boat. It also makes good use of the principle of the "borrowed view," (jiejing) carefully framing views of the surrounding mountains and a famous view of a distant pagoda.
Another existing garden from the Ming dynasty is the Lingering Garden, also in Suzhou, built during the reign of the Wanli Emperor (1573–1620). During the Qing dynasty, twelve tall limestone rocks were added to the garden, symbolizing mountains. The most famous was a picturesque rock called the Auspicious Cloud-Capped Peak, which became a centerpiece of the garden.
A third renowned Ming era garden in Suzhou is the Garden of Cultivation, built during the reign of the Tianqi Emperor (1621–27) by the grandson of Wen Zhengming, a famous Ming painter and calligrapher. The garden is built around a pond, with the Longevity Pavilion on the north side, the Fry Pavilion on the east side, a dramatic rock garden on the south, and the creator's study, the Humble House, to the west.
Qing dynasty (1644–1912)
The Qing dynasty was the last dynasty of China. The most famous gardens in China during this period were the Summer Palace and the Old Summer Palace in Beijing. Both gardens became symbols of luxury and refinement, and were widely described by European visitors.
Father Attiret, a French Jesuit who became court painter for the Qianlong Emperor from 1738 to 1768, described the Jade Terrace of the Isle of Immortality in the Lake of the Summer Palace:
"That which is a true jewel is a rock or island...which is in the middle of this lake, on which is built a small palace, which contains one hundred rooms or salons...of a beauty and a taste which I am not able to express to you. The view is admirable...
Their construction and improvement consumed a large part of the imperial treasury. Empress Dowager Cixi famously diverted money intended for the modernization of the Beiyang Fleet and used it to restore the Summer Palace and the marble teahouse in the shape of boat on Lake Kunming. Both the Summer Palace and Old Summer Palace were destroyed during the Boxer Rebellion and by punitive expeditions of European armies during the nineteenth century, but are now gradually being restored.
In addition to the Old Summer Palace and Summer Palace, between 1703 and 1792 the Qing emperors built a new complex of gardens and palaces in the mountains 200 kilometers northeast of Beijing, to escape the summer heat of the capital. It was called the Chengde Mountain Resort, and it occupied 560 hectares, with seventy-two separate landscape views, recreating landscapes in miniature from many different parts of China. This enormous garden has survived relatively intact.
Renowned scholar gardens which still exist from this period include the Couple's Retreat Garden (1723–1736) and the Retreat & Reflection Garden (1885), both in Suzhou.
Design of the classical garden
A Chinese garden was not meant to be seen all at once; the plan of a classical Chinese garden presented the visitor with a series of perfectly composed and framed glimpses of scenery; a view of a pond, or of a rock, or a grove of bamboo, a blossoming tree, or a view of a distant mountain peak or a pagoda. The 16th-century Chinese writer and philosopher Ji Cheng instructed garden builders to "hide the vulgar and the common as far as the eye can see, and include the excellent and the splendid."
Some early Western visitors to the imperial Chinese gardens felt they were chaotic, crowded with buildings in different styles, without any seeming order. But the Jesuit priest Jean Denis Attiret, who lived in China from 1739 and was a court painter for the Qianlong Emperor, observed there was a "beautiful disorder, an anti-symmetry" in the Chinese garden. "One admires the art with which this irregularity is carried out. Everything is in good taste, and so well arranged, that there is not a single view from which all the beauty can be seen; you have to see it piece by piece."
Chinese classical gardens varied greatly in size. The largest garden in Suzhou, the Humble Administrator's Garden, was a little over ten hectares in area, with one fifth of the garden occupied by the pond. But they did not have to be large. Ji Cheng built a garden for Wu Youyu, the Treasurer of Jinling, that was just under one hectare in size, and the tour of the garden was only four hundred steps long from the entrance to the last viewing point, but Wu Youyu said it contained all the marvels of the province in a single place.
The classical garden was surrounded by a wall, usually painted white, which served as a pure backdrop for the flowers and trees. A pond of water was usually located in the center. Many structures, large and small, were arranged around the pond. In the garden described by Ji Cheng above, the structures occupied two-thirds of the hectare, while the garden itself occupied the other third. In a scholar garden the central building was usually a library or study, connected by galleries with other pavilions which served as observation points of the garden features. These structures also helped divide the garden into individual scenes or landscapes. The other essential elements of a scholar garden were plants, trees, and rocks, all carefully composed into small perfect landscapes. Scholar gardens also often used what was called "borrowed" scenery (借景 jiejing) ; where unexpected views of scenery outside the garden, such as mountain peaks, seemed to be an extension of the garden itself.
Architecture
Chinese gardens are filled with architecture; halls, pavilions, temples, galleries, bridges, kiosks, and towers, occupying a large part of the space. The Humble Administrator's Garden in Suzhou has forty-eight structures, including a residence, several halls for family gatherings and entertainment, eighteen pavilions for viewing different features of the garden, and an assortment of towers, galleries, and bridges, all designed for seeing different parts of the gardens from different points of view. The garden structures are not designed to dominate the landscape, but to be in harmony with it.
Classical gardens traditionally have these structures:
The ceremony hall (ting), or “room”. A building used for family celebrations or ceremonies, usually with an interior courtyard, not far from the entrance gate.
The principal pavilion (da ting), or “large room”, for the reception of guests, for banquets and for celebrating holidays, such as New Years and the Festival of Lanterns. It often has a veranda around the building to provide cool and shade.
The pavilion of flowers (hua ting), or “flower room”. Located near the residence, this building has a rear courtyard filled with flowers, plants, and a small rock garden.
The pavilion facing the four directions (si mian ting), or “four doors room”. This building has folding or movable walls, for opening up a panoramic view of the garden.
The lotus pavilion (he hua ting), or “lotus room”. Built next to a lotus pond, to see the flowers bloom and appreciate their aroma.
The pavilion of mandarin ducks (yuan yang ting), or “mandarin ducks room”. This building is divided into two sections; one facing north used in summer, facing a lotus pond which provided cool air; and the southern part used in winter, with a courtyard planted with pine trees, which remained evergreen, and plum trees, whose blossoms announced the arrival of spring.
In addition to these larger halls and pavilions, the garden is filled with smaller pavilions, (also called ting),or “room”, which are designed for providing shelter from the sun or rain, for contemplating a scene, reciting a poem, taking advantage of a breeze, or simply resting. Pavilions might be located where the dawn can best be watched, where the moonlight shines on the water, where autumn foliage is best seen, where the rain can best be heard on the banana leaves, or where the wind whistles through the bamboo stalks. They are sometimes attached to the wall of another building or sometimes stood by themselves at view points of the garden, by a pond or at the top of a hill. They often are open on three sides.
The names of the pavilions in Chinese gardens express the view or experience they offer the visitor:
The Peak-Worshipping Pavilion (The Lingering Garden) in Suzhou China
The Hall of Distant Fragrances (Humble Administrator's Garden) in Suzhou China
The Mountain View Tower (Humble Administrator's Garden) in Suzhou China
Pavilion of the Moon and Wind (Master of the Nets Garden) in Suzhou China
Pavilion in the Lotus Breeze (Humble Administrator's Garden) in Suzhou China
Listening to the Rain Pavilion (Humble Administrator's Garden) in Suzhou China
Watching the Pines and Appreciating Paintings Hall (Humble Administrator's Garden) in Suzhou China
Spot of Return for Reading (Lingering Garden) in Suzhou China
Between the Mountains and the Water Pavilion (The Couple's Retreat Garden) in Suzhou China
Pavilion Leaning on the Jade (Humble Administrator's Garden) in Suzhou China
Soft Rain Brings Coolness Terrace (Retreat & Reflection Garden) in Suzhou China
Lasting Spring and Moon Viewing Tower (Retreat & Reflection Garden) in Suzhou China
Gardens also often feature two-story towers (lou or ge), usually at the edge of the garden, with a lower story made of stone and a whitewashed upper story, two-thirds the height of the ground floor, which provided a view from above of certain parts of the garden or the distant scenery.
Some gardens have a picturesque stone pavilion in the form of a boat, located in the pond. (called an xie, fang, or shifang). These generally had three parts; a kiosk with winged gables at the front, a more intimate hall in the center, and a two-story structure with a panoramic view of the pond at the rear.
Courtyards (yuan). Gardens contain small enclosed court courtyards, offering quiet and solitude for meditation, painting, drinking tea, or playing on the cithare.
Galleries (lang) are narrow covered corridors which connect the buildings, protect the visitors from the rain and sun, and also help divide the garden into different sections. These galleries are rarely straight; they zigzag or are serpentine, following the wall of the garden, the edge of the pond, or climbing the hill of the rock garden. They have small windows, sometimes round or in odd geometric shapes, to give glimpses of the garden or scenery to those passing through.
Windows and doors are an important architectural feature of the Chinese garden. Sometimes they are round (moon windows or a moon gate) or oval, hexagonal or octagonal, or in the shape of a vase or a piece of fruit. Sometimes they have highly ornamental ceramic frames. The window may carefully frame a branch of a pine tree, or a plum tree in blossom, or another intimate garden scene.
Bridges are another common feature of the Chinese garden. Like the galleries, they are rarely straight, but zigzag (called the Nine-turn bridges) or arch over the ponds, suggesting the bridges of rural China, and providing view points of the garden. Bridges are often built from rough timber or stone-slab raised pathways. Some gardens have brightly painted or lacquered bridges, which give a lighthearted feeling to the garden.
Gardens also often include small, austere houses for solitude and meditation, sometimes in the form of rustic fishing huts, and isolated buildings which serve as libraries or studios (shufang).
Artificial mountains and rock gardens
The artificial mountain (jiashan) or rock garden is an integral element of Chinese classical gardens. The mountain peak was a symbol of virtue, stability and endurance in Confucian philosophy and in the I Ching. A mountain peak on an island was also a central part of the legend of the Isles of the Immortals, and thus became a central element in many classical gardens.
The first rock garden appeared in Chinese garden history in Tu Yuan (literally "Rabbit Garden"), built during the Western Han dynasty (206 BCE – 9 CE). During the Tang dynasty, the rock was elevated to the status of an art object, judged by its form (xing), substance (zhi), color (se), and texture (wen), as well as by its softness, transparency, and other factors. The poet Bo Juyi (772–846) wrote a catalog of the famous rocks of Lake Tai, called Taihu Shiji. These rocks, of limestone sculpted by erosion, became the most highly prized for gardens.
During the Song dynasty, the artificial mountains were made mostly of earth. But Emperor Huizong (1100–1125) nearly ruined the economy of the Song Empire by destroying the bridges of the Grand Canal so he could carry huge rocks by barge to his imperial garden. During the Ming dynasty, the use of piles of rocks to create artificial mountains and grottos reached its peak. During the Qing dynasty, the Ming rock gardens were considered too artificial and the new mountains were composed of both rocks and earth.
The artificial mountain in Chinese gardens today usually has a small view pavilion at the summit. In smaller classical gardens, a single scholar rock represents a mountain, or a row of rocks represents a mountain range.
Water
A pond or lake is the central element of a Chinese garden. The main buildings are usually placed beside it, and pavilions surround the lake to see it from different points of view. The garden usually has a pond for lotus flowers, with a special pavilion for viewing them. There are usually goldfish in the pond, with pavilions over the water for viewing them.
The lake or pond has an important symbolic role in the garden. In the I Ching, water represents lightness and communication, and carried the food of life on its journey through the valleys and plains. It also is the complement to the mountain, the other central element of the garden, and represents dreams and the infinity of spaces. The shape of the garden pond often hides the edges of the pond from viewers on the other side, giving the illusion that the pond goes on to infinity. The softness of the water contrasts with the solidity of the rocks. The water reflects the sky, and therefore is constantly changing, but even a gentle wind can soften or erase the reflections.
The lakes and waterside pavilions in Chinese gardens were also influenced by another classic of Chinese literature, the Shishuo Xinyu by Liu Yiqing (403–444), who described the promenades of the Emperor Jianwen of Jin along the banks of the Hao and the Pu River, in the Garden of the Splendid Forest (Hualin yuan). Many gardens, particularly in the gardens of Jiangnan and the imperial gardens of northern China, have features and names taken from this work.
Small gardens have a single lake, with a rock garden, plants and structures around its edge. Middle-sized gardens will have a single lake with one or more streams coming into the lake, with bridges crossing the streams, or a single long lake divided into two bodies of water by a narrow channel crossed by a bridge. In a very large garden like the Humble Administrator's Garden, the principal feature of the garden is the large lake with its symbolic islands, symbolizing the isles of the immortals. Streams come into the lake, forming additional scenes. Numerous structures give different views of the water, including a stone boat, a covered bridge, and several pavilions by the side of or over the water.
The streams in the Chinese garden always follow a winding course, and are hidden from time to time by rocks or vegetation. A French Jesuit missionary, Father Attiret, who was a painter in the service of the Qianlong Emperor from 1738 to 1768, described one garden he saw:
"The canals are not like those in our country bordered with finely cut stone, but very rustic and lined with pieces or rock, some coming forward, some retreating. which are placed so artistically that you would think it was a work of nature."
Flowers and trees
Flowers and trees, along with water, rocks and architecture, are the fourth essential element of the Chinese garden. They represent nature in its most vivid form, and contrast with the straight lines of the architecture and the permanence, sharp edges and immobility of the rocks. They change continually with the seasons, and provide both sounds (the sound of rain on banana leaves or the wind in the bamboo) and aromas to please the visitor.
Each flower and tree in the garden had its own symbolic meaning. The pine, bamboo and Chinese plum (Prunus mume) were considered the "Three Friends of Winter" (歲寒三友) by the scholars who created classical gardens, prized for remaining green or blooming in winter. They were often painted together by artists like Zhao Mengjian (1199–1264). For scholars, the pine was the emblem of longevity and tenacity, as well as constance in friendship. The bamboo, a hollow straw, represented a wise man, modest and seeking knowledge, and was also noted for being flexible in a storm without breaking. Plum trees were revered as the symbol of rebirth after the winter and the arrival of spring. During the Song dynasty, the favorite tree was the winter plum tree, appreciated for its early pink and white blossoms and sweet aroma.
The peach tree in the Chinese garden symbolized longevity and immortality. Peaches were associated with the classic story The Orchard of Xi Wangmu, the Queen Mother of the West. This story said that in Xi Wangmu's legendary orchard, peach trees flowered only after three thousand years, did not produce fruit for another three thousand years, and did not ripen for another three thousand years. Those who ate these peaches became immortal. This legendary orchard was pictured in many Chinese paintings, and inspired many garden scenes. Pear trees were the symbol of justice and wisdom. The word 'pear' was also a homophone for 'quit' or separate,' and it was considered bad luck to cut a pear, for it would lead to the breakup of a friendship or romance. The pear tree could also symbolize a long friendship or romance, since the tree lived a long time.
The apricot tree symbolized the way of the mandarin, or the government official. During the Tang dynasty, those who passed the imperial examination were rewarded with the banquet in the garden of the apricot trees, or Xingyuan.
The fruit of the pomegranate tree was offered to young couples so they would have male children and numerous descendants. The willow tree represented the friendship and the pleasures of life. Guests were offered willow branches as a symbol of friendship.
Of the flowers in the Chinese garden, the most appreciated were the orchid, peony, and lotus (Nelumbo nucifera). During the Tang dynasty, the peony, the symbol of opulence and a flower with a delicate fragrance, was the most celebrated flower in the garden. The poet Zhou Dunyi wrote a famous elegy to the lotus, comparing it to a junzi, a man who possessed integrity and balance. The orchid was the symbol of nobility, and of impossible love, as in the Chinese expression "a faraway orchid in a lonely valley." The lotus was admired for its purity, and its efforts to reach out of the water to flower in the air made it a symbol of the search for knowledge. The chrysanthemum was elegized the poet Tao Yuanming, who surrounded his hermit's hut with the flower, and wrote a famous verse:
"At the feet of the Eastern fence, I pick a chrysanthemum,
In the distance, detached and serene, I see the Mountains of the South."
The creators of the Chinese garden were careful to preserve the natural appearance of the landscape. Trimming and root pruning, if done at all, tried to preserve the natural form. Dwarf trees that were gnarled and ancient-looking were particularly prized in the miniature landscapes of Chinese gardens.
"Borrowing scenery", time and seasons
According to Ji Cheng's 16th century book Yuanye, "The Craft of Gardens," "borrowed scenery" (jiejing) was the most important thing of a garden. This could mean using scenes outside the garden, such as a view of distant mountains or the trees in the neighboring garden, to create the illusion that garden was much bigger than it was. The most famous example was the mist-shrouded view of the North Temple Pagoda in Suzhou, seen in the distance over the pond of the Humble Administrator's Garden.
But, as Ji Cheng wrote, it could also be "the immaculate ribbon of a stream, animals, birds, fish, or other natural elements (rain, wind, snow), or something less tangible, such as a moonbeam, a reflection in a lake, morning mist, or the red sky of a sunset." It could also be a sound; he recommended locating a pavilion near a temple, so that the chanted prayers could be heard; planting fragrant flowers next to paths and pavilions, so visitors would appreciate their aromas; that bird perches be created to encourage birds to come to sing in the garden, that streams be designed to make pleasant sounds, and that banana trees be planted in courtyards so the rain would patter on their leaves. "A judicious 'borrowing' does not have a reason." Ji Cheng wrote. "It is born simply of feeling created by the beauty of a scene."
The season and the time of day were also important elements. Garden designers took into account the scenes of the garden that would look best in winter, summer, spring and autumn, and those best viewed at night, in the morning or afternoon. Ji Cheng wrote: "In the heart of the tumult of the city, you should choose visions that are serene and refined: from a raised clearing, you look to the distant horizon, surrounded by mountains like a screen; in an open pavilion, a gentle and light breeze invades the room; from the front door, the running water of spring flows toward the marsh."
Actually borrowing scenery is the conclusive, last chapter of Yuanye that explains borrowing scenery as a holistic understanding of the essence of landscape design in its entirety. The ever-changing moods and appearances of nature in a given landscape in full action are understood by the author as an independent function that becomes an agent for garden making. It is nature including the garden maker that creates.
Concealment and surprise
Another important garden element was concealment and surprise. The garden was not meant to be seen all at once, it was laid out to present a series of scenes. Visitors moved from scene to scene either within enclosed galleries or by winding paths which concealed the scenes until the last moment. The scenes would suddenly appear at the turn of a path, through a window, or hidden behind a screen of bamboo. They might be revealed through round "moon doors" or through windows of unusual shapes, or windows with elaborate lattices that broke the view into pieces.
In art and literature
The garden plays an important part in Chinese art and literature, and at the same time art and literature have inspired many gardens. The school of painting called "Shanshui" (literally 'mountains and water' and with the actual meaning of 'landscape'), which began in the 5th century, established the principles of Chinese landscape painting, which were very similar to those of Chinese gardening. These paintings were not meant to be realistic; they were meant to portray what the artist felt, rather than what he saw.
The landscape painter Shitao (1641–1720) wrote that he wanted to "'...create a landscape which was not spoiled by any vulgar banality..." He wanted to create a sense of vertigo in the viewer: "to express a universe inaccessible to man, without any route that led there, like the isles of Bohai, Penglan and Fanghu, where only the immortals can live, and which a man cannot imagine. That is the vertigo that exists in the natural universe. To express it in painting, you must show jagged peaks, precipices, hanging bridges, great chasms. For the effect to be truly marvelous, it must be done purely by the force of the brush." This was the emotion that garden designers wanted to create with their scholar rocks and miniature mountain ranges.
In his book, Craft of Gardens, the garden designer Ji Cheng wrote: "The spirit and the charm of mountains and forests must be studied in depth; ...only the knowledge of the real permits the creation of the artificial, so that the work created possesses the spirit of the real, in part because of divine inspiration, but especially because of human effort." He described the effect he wanted to achieve in the design of an autumn garden scene: "The feelings are in harmony with the purity, with the sense of withdrawal. The spirit rejoices at the mountains and ravines. Suddenly the spirit, detached from the world of small things, is animated and seems to penetrate to the interior of a painting, and to promenade there..."
In literature, gardens were frequently the subject of the genre of poetry called "Tianyuan", literally 'fields and gardens,' which reached its peak in the Tang dynasty (618–907) with such poets as Wang Wei (701–761). The names of the Surging Waves Garden and the Garden of Meditation in Suzhou are taken from lines of Chinese poetry. Within the gardens, the individual pavilions and view points were frequently dedicated to verses of poems, inscribed on stones or plaques. The Moon Comes with the Breeze Pavilion at the Couple's Retreat Garden, used for moon-viewing, has the inscription of a verse by Han Yu:
"The twilight brings the Autumn
And the wind brings the moon here."
And the Peony Hall in the Couple's Retreat Garden is dedicated to a verse by Li Bai:
"The spring breeze is gently stroking the balustrade
and the peony is wet with dew."
Wang Wei (701–761) was a poet, painter and Buddhist monk, who worked first as a court official before retiring to Lantian, where he built one of the first wenren yuan, or scholar's gardens, called the Valley of the Jante. In this garden, a series of twenty scenes, like the paintings of a scroll or album, unrolled before the viewer, each illustrated by a verse of poetry. For example, one scene illustrated this poem:
"The white rock emerges from the torrent;
The cold sky with red leaves scattering:
On the mountain path, the rain is fleeing,
the blue of the emptiness dampens our clothes."
The Valley of the Jante garden disappeared, but its memory, preserved in paintings and poems, inspired many other scholar's gardens.
The social and cultural importance of the garden is illustrated in the classical novel Dream of the Red Chamber by Cao Xueqin which unfolds almost exclusively in a garden.
Philosophy
The Chinese classical garden had multiple functions. It could be used for banquets, celebrations, reunions, or romance. It could be used to find solitude and for contemplation. It was a calm place for painting, poetry, calligraphy, and music, and for studying classic texts. It was a place for drinking tea and for poets to become happily drunk on wine. It was a showcase to display the cultivation and aesthetic taste of the owner. But it also had a philosophical message.
Taoism had a strong influence on the classical garden. After the Han dynasty (206 BC – 220 AD), gardens were frequently constructed as retreats for government officials who had lost their posts or who wanted to escape the pressures and corruption of court life in the capital. They chose to pursue the Taoist ideals of disengagement from worldly concerns.
For followers of Taoism, enlightenment could be reached by contemplation of the unity of creation, in which order and harmony are inherent to the natural world.
The gardens were intended to evoke the idyllic feeling of wandering through a natural landscape, to feel closer to the ancient way of life, and to appreciate the harmony between man and nature.
In Taoism, rocks and water were opposites, yin and yang, but they complemented and completed one another. Rocks were solid but water could wear away rock. The deeply eroded rocks from Lake Tai used in the classical garden illustrated this principle.
Borrowing scenery is a most fundamental idea in Ming period garden making theory (see above).
The winding paths and zig-zag galleries bridges that led visitors from one garden scene to another also had a message. They illustrated a Chinese proverb, "By detours, access to secrets".
According to the landscape historian and architect Che Bing Chiu, every garden was "a quest for paradise. of a lost world, of a utopian universe. The scholar's garden participated in this quest; on the one hand the quest for the home of the Immortals, on the other hand the search for the world of the golden age so dear to the heart of the scholar."
A more recent view of the philosophy of the garden was expressed by Zhou Ganzhi, the President of the Chinese Society of Landscape Architecture, and Academician at the Chinese Academy of Sciences and the Chinese Academy of Engineering, in 2007: "Chinese classical gardens are a perfect integration of nature and work by man. They are an imitation of nature, and fully manifest the beauty of nature. They can also be seen as an improvement on nature; one from which the light of human artistic genius shines."
Influence
Chinese influence on the Japanese garden
The Chinese classical garden had a notable influence on the early Japanese garden. The influence of China first reached Japan through Korea before 600 AD. In 607 AD, the Japanese crown prince Shotoku sent a diplomatic mission to the Chinese court, which began a cultural exchange lasting for centuries. Hundreds of Japanese scholars were sent to study the Chinese language, political system, and culture. The Japanese Ambassador to China, Ono no Imoko, described the great landscape gardens of the Chinese Emperor to the Japanese court. His reports had a profound influence on the development of Japanese landscape design.
During the Nara period (710-794), when the Japanese capital was located at Nara, and later at Heian, the Japanese court created large landscape gardens with lakes and pavilions on the Chinese model for aristocrats to promenade and to drift leisurely in small boats, and more intimate gardens for contemplation and religious meditation.
A Japanese monk named Eisai (1141–1215) imported the Rinzai school of Zen Buddhism from China to Japan, which led to the creation of a famous and unique Japanese gardening style, the Zen garden, exemplified by the garden of Ryōan-ji. He also brought green tea from China to Japan, originally to keep monks awake during long meditation, giving the basis for the Japanese tea ceremony, which became an important ritual in Japanese gardens.
The Japanese garden designer Muso Soseki (1275–1351) created the celebrated Moss Garden (Kokedera) in Kyoto, which included a recreation of the Isles of Eight Immortals, called Horai in Japanese, which were an important feature of many Chinese gardens. During the Kamakura period (1185–1333), and particularly during the Muromachi period (1336–1573) the Japanese garden became more austere than the Chinese garden, following its own aesthetic principles.
In Europe
The first European to describe a Chinese garden was the Venetian merchant and traveler Marco Polo, who visited the summer palace of Kublai Khan at Xanadu. The garden of Kublai Khan had a later effect on European culture; In 1797, it inspired the romantic poem, Kubla Khan, by the English romantic poet Samuel Taylor Coleridge.
Marco Polo also described the gardens of the imperial palace in Khanbaliq, the Mongol name for the city which eventually became Beijing. He described ramparts, balustrades and pavilions surrounding a deep lake full of fish and with swans and other aquatic birds; whose central feature was a manmade hill one hundred steps high and a thousand steps around, covered with evergreen trees and decorated with green azurite stones.
The first Jesuit priest, Francis Xavier, arrived in China in 1552, and the priest Matteo Ricci received permission to settle in Beijing in 1601. Jesuit priests began sending accounts of Chinese culture and gardens to Europe. Louis Le Comte, the mathematician to the King of France, travelled to China in 1685. He described how the Chinese gardens had grottos, artificial hills and rocks piled to imitate nature, and did not arrange their gardens geometrically.
In the 18th century, as Chinese vases and other decorative objects began to arrive in Europe, there was a surge of popularity for Chinoiserie. The painters Watteau and François Boucher painted Chinese scenes as they imagined them, and Catherine the Great decorated a room in her palace in Chinese style. There was great interest in everything Chinese, including gardens.
In 1738, the French Jesuit missionary and painter Jean Denis Attiret, went to China, where he became court painter to the Qianlong Emperor. He described in great detail what he saw in the imperial gardens near Beijing:
"One comes out of a valley, not by a straight wide alley as in Europe, but by zigzags, by roundabout paths, each one ornamented with small pavilions and grottos, and when you exit one valley you find yourself in another, different from the first in the form of the landscape or the style of the buildings. All the mountains and hills are covered with flowering trees, which are very common here. It is a true terrestrial paradise. The canals are not at all like ours- bordered with cut stone- they are rustic, with pieces of rock, some leaning forward, some backwards, placed with such art you would think they were natural. Sometimes a canal is wide, sometimes narrow. Here they twist, there they curve, as if they were really created by the hills and rocks. The edges are planted with flowers in rock gardens, which seem to have been created by nature. Each season has its own flowers. Aside from the canals, everywhere there are paths paved with small stones, which lead from one valley to the other. These paths also twist and turn, sometimes coming close to the canals, sometimes far away."
Attiret wrote:
"Everything is truly great and beautiful, both as to the design and the execution: and [the gardens] struck me the more, because I had never seen any thing that bore any manner of resemblance to them, in any part of the world that I had been before."
The Qianlong Emperor (1711–1799) was equally interested in what was going on in Europe. He commissioned the Jesuit priest Father Castiglione, who was trained in engineering, to build fountains for his garden similar to those he had heard about in the gardens at Versailles.
Chinese architecture and aesthetics may also have influenced the English landscape garden style. In 1685, the English diplomat and writer Sir William Temple wrote an essay Upon the garden of Epicurus (published in 1692), a passage in which contrasted European theories of symmetrical gardens with asymmetrical compositions from China. Temple had never visited China, but had heard of Chinese (or Japanese) gardens, perhaps in the Netherlands. He noted that Chinese gardens avoided formal rows of trees and flower beds, and instead placed trees, plants, and other garden features in irregular ways to strike the eye and create beautiful compositions. He gave the term Sharawadgi to this approach. His observations on the Chinese garden were cited by the essayist Joseph Addison in an essay in 1712, who used them to attack the English gardeners who, instead of imitating nature, tried to make their gardens in the French style, as far from nature as possible.
The English landscape garden was already well-established in England in the first part of the 18th century, influenced by the travel to Italy by the British upper class and their desire to have a new style of garden to match the Palladian style of architecture they chose for their country houses, and by the romantic landscapes of Claude Lorraine and other painters, but the novelty and exoticism of Chinese art and architecture in Europe led in 1738 to the construction of the first Chinese house in an English garden, at Stowe House, alongside Roman temples, Gothic ruins and other architectural styles.
The style became even more popular thanks to William Chambers (1723–1796), who lived in China from 1745 to 1747, and wrote a book, The Drawings, buildings, furniture, habits, machines and untensils of the Chinese, published in 1757. He urged western garden designers to use Chinese stylistic conventions such as concealment, asymmetry, and naturalism. Later, in 1772, Chambers published his Dissertation on Oriental Gardening, a rather fanciful elaboration of contemporary ideas about the naturalistic style of gardening in China.
Chambers was a fierce critic of Capability Brown, the leading designer of the English landscape garden, which Chambers considered boring. Chambers believed that gardens should be full of surprises. In 1761 he built the Great Pagoda in Kew Gardens, London, along with a mosque, a temple of the sun, a ruined arch, and Palladian bridge. Thanks to Chambers Chinese structures began to appear in other English gardens, then in France and elsewhere on the continent. Carmontelle added a Chinese pavilion to his garden at Parc Monceau in Paris (1772), and the Duc de Choiseul built a pagoda on his estate at Chanteloup between 1775 and 1778, now the only part of the estate to survive. The Russian Empress Catherine the Great built her own pagoda in the garden of her palace of Tsarskoye Selo, near Saint Petersburg, between 1778 and 1786. Many French critics disliked the term "English Garden", so they began to use the term 'Anglo-Chinois" to describe the style.
| Technology | Buildings and infrastructure | null |
1590904 | https://en.wikipedia.org/wiki/Precipitation%20hardening | Precipitation hardening | Precipitation hardening, also called age hardening or particle hardening, is a heat treatment technique used to increase the yield strength of malleable materials, including most structural alloys of aluminium, magnesium, nickel, titanium, and some steels, stainless steels, and duplex stainless steel. In superalloys, it is known to cause yield strength anomaly providing excellent high-temperature strength.
Precipitation hardening relies on changes in solid solubility with temperature to produce fine particles of an impurity phase, which impede the movement of dislocations, or defects in a crystal's lattice. Since dislocations are often the dominant carriers of plasticity, this serves to harden the material. The impurities play the same role as the particle substances in particle-reinforced composite materials. Just as the formation of ice in air can produce clouds, snow, or hail, depending upon the thermal history of a given portion of the atmosphere, precipitation in solids can produce many different sizes of particles, which have radically different properties. Unlike ordinary tempering, alloys must be kept at elevated temperature for hours to allow precipitation to take place. This time delay is called "aging". Solution treatment and aging is sometimes abbreviated "STA" in specifications and certificates for metals.
Two different heat treatments involving precipitates can alter the strength of a material: solution heat treating and precipitation heat treating. Solid solution strengthening involves formation of a single-phase solid solution via quenching. Precipitation heat treating involves the addition of impurity particles to increase a material's strength.
Kinetics versus thermodynamics
This technique exploits the phenomenon of supersaturation, and involves careful balancing of the driving force for precipitation and the thermal activation energy available for both desirable and undesirable processes.
Nucleation occurs at a relatively high temperature (often just below the solubility limit) so that the kinetic barrier of surface energy can be more easily overcome and the maximum number of precipitate particles can form. These particles are then allowed to grow at lower temperature in a process called ageing. This is carried out under conditions of low solubility so that thermodynamics drive a greater total volume of precipitate formation.
Diffusion's exponential dependence upon temperature makes precipitation strengthening, like all heat treatments, a fairly delicate process. Too little diffusion (under ageing), and the particles will be too small to impede dislocations effectively; too much (over ageing), and they will be too large and dispersed to interact with the majority of dislocations.
Alloy design
Precipitation strengthening is possible if the line of solid solubility slopes strongly toward the center of a phase diagram. While a large volume of precipitate particles is desirable, a small enough amount of the alloying element should be added so that it remains easily soluble at some reasonable annealing temperature. Although large volumes are often wanted, they are wanted in small particle sizes as to avoid a decrease in strength as is explained below.
Elements used for precipitation strengthening in typical aluminium and titanium alloys make up about 10% of their composition. While binary alloys are more easily understood as an academic exercise, commercial alloys often use three components for precipitation strengthening, in compositions such as Al(Mg, Cu) and Ti(Al, V). A large number of other constituents may be unintentional, but benign, or may be added for other purposes such as grain refinement or corrosion resistance. An example is the addition of Sc and Zr to aluminum alloys to form FCC L12 structures that help refine grains and strengthen the material. In some cases, such as many aluminium alloys, an increase in strength is achieved at the expense of corrosion resistance. More recent technology is focused on additive manufacturing due to the higher amount of metastable phases that can be obtained due to the fast cooling, whereas traditional casting is more limited to equilibrium phases.
The addition of large amounts of nickel and chromium needed for corrosion resistance in stainless steels means that traditional hardening and tempering methods are not effective. However, precipitates of chromium, copper, or other elements can strengthen the steel by similar amounts in comparison to hardening and tempering. The strength can be tailored by adjusting the annealing process, with lower initial temperatures resulting in higher strengths. The lower initial temperatures increase the driving force of nucleation. More driving force means more nucleation sites, and more sites means more places for dislocations to be disrupted while the finished part is in use.
Many alloy systems allow the ageing temperature to be adjusted. For instance, some aluminium alloys used to make rivets for aircraft construction are kept in dry ice from their initial heat treatment until they are installed in the structure. After this type of rivet is deformed into its final shape, ageing occurs at room temperature and increases its strength, locking the structure together. Higher ageing temperatures would risk over-ageing other parts of the structure, and require expensive post-assembly heat treatment because a high ageing temperature promotes the precipitate to grow too readily.
Types of hardening
There are several ways by which a matrix can be hardened by precipitates, which could also be different for deforming precipitates and non-deforming precipitates.
Deforming particles (weak precipitates):
Coherency hardening occurs when the interface between the particles and the matrix is coherent, which depends on parameters like particle size and the way that particles are introduced. Coherency is where the lattice of the precipitate and that of the matrix are continuous across the interface. Small particles precipitated from supersaturated solid solution usually have coherent interfaces with the matrix. Coherency hardening originates from the atomic volume difference between precipitate and the matrix, which results in a coherency strain. If the atomic volume of the precipitate is smaller, there will be tension because the lattice atoms are located closer than their normal conditions while when the atomic volume of the precipitate is larger, there will be compression of the lattice atoms, as they are further apart than their normal position. Regardless of whether the lattice is under compression or tension, the associated stress field interacts with dislocations leading to decreased dislocation motion either by repulsion or attraction of the dislocations, leading to an increase in yield strength, similar to the size effect in solid solution strengthening. What differentiates this mechanism from solid solution strengthening is the fact that the precipitate has a definite size, not an atom, and therefore a stronger interaction with dislocations.
Modulus hardening results from the different shear modulus of the precipitate and the matrix, which leads to an energy change of dislocation line tension when the dislocation line cuts the precipitate. Also, the dislocation line could bend when entering the precipitate, increasing the affected length of the dislocation line. Again, the strengthening arises in a way similar to that of solid solution strengthening, where there is a mismatch in the lattice that interacts with the dislocations, impeding their motion. Of course, the severity of the interaction is different than that of solid solution and coherency strengthening.
Chemical strengthening is associated with the surface energy of the newly introduced precipitate-matrix interface when the particle is sheared by dislocations. Because it takes energy to make the surface, some of the stress that is causing dislocation motion is accommodated by the additional surfaces. Like modulus hardening, the analysis of interfacial area can be complicated by dislocation line distortion.
Order strengthening occurs when the precipitate is an ordered structure such that bond energy before and after shearing is different. For example, in an ordered cubic crystal with composition AB, the bond energy of A-A and B-B after shearing is higher than that of the A-B bond before. The associated energy increase per unit area is anti-phase boundary energy and accumulates gradually as the dislocation passes through the particle. However, a second dislocation could remove the anti-phase domain left by the first dislocation when traverses the particle. The attraction of the particle and the repulsion of the first dislocation maintains a balanced distance between two dislocations, which makes order strengthening more complicated. Except for when there are very fine particles, this mechanism is generally not as effective as others to strengthen. Another way to consider this mechanism is that when a dislocation shears a particle, the stacking sequence between the new surface made and the matrix is broken, and the bonding is not stable. To get the sequence back into this interface, another dislocation, is needed to shift the stacking. The first and second dislocation are often called a superdislocation. Because superdislocations are required to shear these particles, there is strengthening because of the decreased dislocation motion.
Non-deforming particles (strong precipitate):
In non-deforming particles, where the spacing is small enough or the precipitate-matrix interface is disordered, dislocation bows instead of shears. The strengthening is related to the effective spacing between particles considering finite particle size, but not particle strength, because once the particle is strong enough for the dislocations to bow rather than cut, further increase of the dislocation penetration resistance won't affect strengthening. The main mechanism therefore is Orowan strengthening, where the strong particles do not allow for dislocations to move past. Therefore bowing must occur and in this bowing can cause dislocation loops to build up, which decreases the space available for additional dislocation to bow between. If the dislocations cannot shear particles and cannot move past them, then dislocation motion is successfully impeded.
Theory
The primary species of precipitation strengthening are second phase particles. These particles impede the movement of dislocations throughout the lattice. You can determine whether or not second phase particles will precipitate into solution from the solidus line on the phase diagram for the particles. Physically, this strengthening effect can be attributed both to size and modulus effects, and to interfacial or surface energy.
The presence of second phase particles often causes lattice distortions. These lattice distortions result when the precipitate particles differ in size and crystallographic structure from the host atoms. Smaller precipitate particles in a host lattice leads to a tensile stress, whereas larger precipitate particles leads to a compressive stress. Dislocation defects also create a stress field. Above the dislocation there is a compressive stress and below there is a tensile stress. Consequently, there is a negative interaction energy between a dislocation and a precipitate that each respectively cause a compressive and a tensile stress or vice versa. In other words, the dislocation will be attracted to the precipitate. In addition, there is a positive interaction energy between a dislocation and a precipitate that have the same type of stress field. This means that the dislocation will be repulsed by the precipitate.
Precipitate particles also serve by locally changing the stiffness of a material. Dislocations are repulsed by regions of higher stiffness. Conversely, if the precipitate causes the material to be locally more compliant, then the dislocation will be attracted to that region. In addition, there are three types of interphase boundaries (IPBs).
The first type is a coherent or ordered IPB, the atoms match up one by one along the boundary. Due to the difference in lattice parameters of the two phases, a coherency strain energy is associated with this type of boundary. The second type is a fully disordered IPB and there are no coherency strains, but the particle tends to be non-deforming to dislocations. The last one is a partially ordered IPB, so coherency strains are partially relieved by the periodic introduction of dislocations along the boundary.
In coherent precipitates in a matrix, if the precipitate has a lattice parameter less than that of the matrix, then the atomic match across the IPB leads to an internal stress field that interacts with moving dislocations.
There are two deformation paths, one is the coherency hardening, the lattice mismatch is
Where is the shear modulus, is the coherent lattice mismatch, is the particle radius, is the particle volume fraction, is the burgers vector, equals the concentration.
The other one is modulus hardening. The energy of the dislocation energy is , when it cuts through the precipitate, its energy is , the change in line segment energy is
.
The maximum dislocation length affected is the particle diameter, the line tension change takes place gradually over a distance equal to . The interaction force between the dislocation and the precipitate is
and .
Furthermore, a dislocation may cut through a precipitate particle, and introduce more precipitate-matrix interface, which is chemical strengthening. When the dislocation is entering the particle and is within the particle, the upper part of the particle shears b with respect to the lower part accompanies the dislocation entry. A similar process occurs when the dislocation exits the particle. The complete transit is accompanied by creation of matrix-precipitate surface area of approximate magnitude , where r is the radius of the particle and b is the magnitude of the burgers vector. The resulting increase in surface energy is , where is the surface energy. The maximum force between the dislocation and particle is , the corresponding flow stress should be .
When a particle is sheared by a dislocation, a threshold shear stress is needed to deform the particle. The expression for the required shear stress is as follows:
When the precipitate size is small, the required shear stress is proportional to the precipitate size , However, for a fixed particle volume fraction, this stress may decrease at larger values of r owing to an increase in particle spacing. The overall level of the curve is raised by increases in either inherent particle strength or particle volume fraction.
The dislocation can also bow around a precipitate particle through so-called Orowan mechanism.
Since the particle is non-deforming, the dislocation bows around the particles (), the stress required to effect the bypassing is inversely proportional to the interparticle spacing , that is, , where is the particle radius. Dislocation loops encircle the particles after the bypass operation, a subsequent dislocation would have to be extruded between the loops. Thus, the effective particle spacing for the second dislocation is reduced to with , and the bypassing stress for this dislocation should be , which is greater than for the first one. However, as the radius of particle increases, will increase so as to maintain the same volume fraction of precipitates, will increase and will decrease. As a result, the material will become weaker as the precipitate size increases.
For a fixed particle volume fraction, decreases with increasing r as this is accompanied by an increase in particle spacing.
On the other hand, increasing increases the level of the stress as a result of a finer particle spacing. The level of is unaffected by particle strength. That is, once a particle is strong enough to resist cutting, any further increase in its resistance to dislocation penetration has no effect on , which depends only on matrix properties and effective particle spacing.
If particles of A of volume fraction are dispersed in a matrix, particles are sheared for and are bypassed for , maximum strength is obtained at , where the cutting and bowing stresses are equal. If inherently harder particles of B of the same volume fraction are present, the level of the curve is increased but that of the one is not. Maximum hardening, greater than that for A particles, is found at . Increasing the volume fraction of A raises the level of both and and increases the maximum strength obtained. The latter is found at , which may be either less than or greater than depending on the shape of the curve.
Governing equations
There are two main types of equations to describe the two mechanisms for precipitation hardening based on weak and strong precipitates. Weak precipitates can be sheared by dislocations while strong precipitates cannot, and therefore the dislocation must bow. First, it is important to consider the difference between these two different mechanisms in terms of the dislocation line tension that they make. The line tension balance equation is:
Where is the radius of the dislocation at a certain stress. Strong obstacles have small due to the bowing of the dislocation. Still, decreasing obstacle strength will increase the and must be included in the calculation. L’ is also equal to the effective spacing between obstacles L. This leaves an equation for strong obstacles:
Considering weak particles, should be nearing due to the dislocation line staying relatively straight through obstacles. Also , L’ will be:
which states the weak particle equation:
Now, consider the mechanisms for each regime:
Dislocation cutting through particles:
For most strengthening at the early stage, it increases with , where is a dimensionless mismatch parameter (for example, in coherency hardening, is the fractional change of precipitate and matrix lattice parameter), is the volume fraction of precipitate, is the precipitate radius, and is the magnitude of the Burgers vector. According to this relationship, materials strength increases with increasing mismatch, volume fraction, and particle size, so that dislocation is easier to cut through particles with smaller radius.
For different types of hardening through cutting, governing equations are as following.
For coherency hardening,
,
,
where is increased shear stress, is the shear modulus of the matrix, and are the lattice parameter of the precipitate or the matrix.
For modulus hardening,
,
,
where and are the shear modulus of the precipitate or the matrix.
For chemical strengthening,
,
,
where is the particle-matrix interphase surface energy.
For order strengthening,
(low , early stage precipitation), where the dislocations are widely separated;
(high , early stage precipitation), where the dislocations are not widely separated;
, where is anti-phase boundary energy.
Dislocations bowing around particles: When the precipitate is strong enough to resist dislocation penetration, dislocation bows and the maximum stress is given by the Orowan equation. Dislocation bowing, also called Orowan strengthening, is more likely to occur when the particle density in the material is lower.
where is the material strength, is the shear modulus, is the magnitude of the Burgers vector, is the distance between pinning points, and is the second phase particle radius. This governing equation shows that for dislocation bowing the strength is inversely proportional to the second phase particle radius , because when the volume fraction of the precipitate is fixed, the spacing between particles increases concurrently with the particle radius , therefore increases with .
These governing equations show that the precipitation hardening mechanism depends on the size of the precipitate particles. At small , cutting will dominate, while at large , bowing will dominate.
Looking at the plot of both equations, it is clear that there is a critical radius at which max strengthening occurs. This critical radius is typically 5-30 nm.
The Orowan strengthening model above neglects changes to the dislocations due to the bending. If bowing is accounted for, and the instability condition in the Frank-Read mechanism is assumed, the critical stress for dislocations bowing between pinning segments can be described as:
where is a function of , is the angle between the dislocation line and the Burgers vector, is the effective particle separation, is the Burgers vector, and is the particle radius.
Other Considerations
Grain Size Control
Precipitates in a polycrystalline material can act as grain refiners if they are nucleated or located near grain boundaries, where they pin the grain boundaries as an alloy is solidifying and do not allow for a coarse microstructure. This is helpful, as finer microstructures often outperform (mechanical properties) coarser ones at room temperatures. In recent times nano-precipitates are being studied under creep conditions. These precipitates can also pin the grain boundary at higher temperatures, essentially acting as "friction". Another useful effect can be to impede grain-boundary sliding under diffusional creep conditions with very fine precipitates and if the precipitates are homogeneously dispersed in the matrix, then these same precipitates in the grains might interact with dislocations under creep dislocation creep conditions.
Secondary Precipitates
Different precipitates, depending on their elemental compositions, can form under certain aging conditions that were not previously there. Secondary precipitates can arise from removing solutes from the matrix solid solution states. The control of this can be exploited to control the microstructure and influence properties.
Computational discovery of new alloys
While significant effort has been made to develop new alloys, the experimental results take time and money to be implemented. One possible alternative is doing simulations with Density functional theory, that can take advantage of, in the context of precipitation hardening, the crystalline structure of precipitates and of the matrix and allow the exploration of a lot more alternatives than with experiments in the traditional form.
One strategy for doing these simulations is focusing on the ordered structures that can be found in many metal alloys, like the long-period stacking ordered (LPSO) structures that have been observed in numerous systems. The LPSO structure is long packed layered configuration along one axis with some layers enriched with precipitated elements. This allows to exploit the symmetry of the supercells and it suits well with the currently available DFT methods.
In this way, some researchers have developed strategies to screen the possible strengthening precipitates that allow decreasing the weight of some metal alloys. For example, Mg-alloys have received progressive interest to replace Aluminum and Steel in the vehicle industry because it is one of the lighter structural metals. However, Mg-alloys show issues with low strength and ductility which have limited their use. To overcome this, the Precipitation hardening technique, through the addition of rare earth elements, has been used to improve the alloy strength and ductility. Specifically, the LPSO structures were found that are responsible for these increments, generating an Mg-alloy that exhibited high-yield strength: 610 MPa at 5% of elongation at room temperature.
In this way, some researchers have developed strategies to Looking for cheaper alternatives than Rare Elements (RE) it was simulated a ternary system with Mg-Xl-Xs, where Xl and Xs correspond to atoms larger than and shorter than Mg, respectively. Under this study, it was confirmed more than 85 Mg-Re-Xs LPSO structures, showing the DFT ability to predict known LPSO ternary structures. Then they explore the 11 non-RE Xl elements and was found that 4 of them are thermodynamically stable. One of them is the Mg-Ca-Zn system that is predicted to form an LPSO structure.
Following the previous DFT predictions, other investigators made experiments with the Mg-Zn-Y-Mn-Ca system and found that at 0.34%at Ca addition the mechanical properties of the system were enhanced due to the formation of LPSO-structures, achieving “a good balance of the strength and ductibility”.
Examples of precipitation hardening materials
2000-series aluminium alloys (important examples: 2024 and 2019, also Y alloy and Hiduminium)
6000-series aluminium alloys (important example: 6061 for bicycle frames and aeronautical structures)
7000-series aluminium alloys (important examples: 7075 and 7475)
17-4 stainless steel (UNS S17400)
Maraging steel
Inconel 718
Alloy X-750
René 41
Waspaloy
Mulberry (uranium alloy)
NAK55 Low Carbon Steel
| Technology | Metallurgy | null |
1591064 | https://en.wikipedia.org/wiki/Breathalyzer | Breathalyzer | A breathalyzer or breathalyser (a portmanteau of breath and analyzer/analyser), also called an alcohol meter, is a device for measuring breath alcohol content (BrAC). It is commonly utilized by law enforcement officers whenever they initiate traffic stops. The name is a genericized trademark of the Breathalyzer brand name of instruments developed by inventor Patrick Tegeler in the 1950s.
Origins
Research into the possibilities of using breath to test for alcohol in a person's body dates as far back as 1874, when Francis E. Anstie made the observation that small amounts of alcohol were excreted in breath.
In 1927, Emil Bogen produced a paper on breath analysis. He collected air in a football bladder and then tested this air for traces of alcohol, discovering that the alcohol content of 2 litres of expired air was a little greater than that of 1 cc of urine. Also in 1927, a Chicago chemist, William Duncan McNally, invented a breathalyzer in which the breath moving through chemicals in water would change color. One suggested use for his invention was for housewives to test whether their husbands had been drinking. In December 1927, in a case in Marlborough, England, Dr. Gorsky, a police surgeon, asked a suspect to inflate a football bladder with his breath. Since the 2 liters of the man's breath contained 1.5 mg of ethanol, Gorsky testified before the court that the defendant was "50% drunk". The use of drunkenness as the standard, as opposed to BAC, perhaps invalidated the analysis, as tolerance to alcohol varies. However, the story illustrates the general principles of breath analysis.
In 1931 the first practical roadside breath-testing device was the drunkometer developed by Rolla Neil Harger of the Indiana University School of Medicine. The drunkometer collected a motorist's breath sample directly into a balloon inside the machine. The breath sample was then pumped through an acidified potassium permanganate solution. If there was alcohol in the breath sample, the solution changed color. The greater the color change, the more alcohol there was present in the breath. The drunkometer was manufactured and sold by Stephenson Corporation of Red Bank, New Jersey.
In 1954 Robert Frank Borkenstein (1912–2002) was a captain with the Indiana State Police and later a professor at Indiana University Bloomington. His trademarked Breathalyzer used chemical oxidation and photometry to determine alcohol concentrations. The invention of the Breathalyzer provided law enforcement with a quick and portable test to determine an individual's intoxication level via breath analysis.
Subsequent breath analyzers have converted primarily to infrared spectroscopy. In 1967 in Britain, Bill Ducie and Tom Parry Jones developed and marketed the first electronic breathalyser. They established Lion Laboratories in Cardiff. Ducie was a chartered electrical engineer, and Tom Parry Jones was a lecturer at UWIST. The Road Safety Act 1967 introduced the first legally enforceable maximum blood alcohol level for drivers in the UK, above which it became an offence to be in charge of a motor vehicle; and introduced the roadside breathalyser, made available to police forces across the country. In 1979, Lion Laboratories' version of the breathalyser, known as the Alcolyser and incorporating crystal-filled tubes that changed colour above a certain level of alcohol in the breath, was approved for police use. Lion Laboratories won the Queen's Award for Technological Achievement for the product in 1980, and it began to be marketed worldwide. The Alcolyser was superseded by the Lion Intoximeter 3000 in 1983, and later by the Lion Alcolmeter and Lion Intoxilyser. These later models used a fuel cell alcohol sensor rather than crystals, providing a more reliable curbside test and removing the need for blood or urine samples to be taken at a police station. In 1991, Lion Laboratories was sold to the American company MPD, Inc.
Accuracy
Breath analyzers do not directly measure blood alcohol concentration (BAC), which requires the analysis of a blood sample. Instead, they measure the amount of alcohol in one's breath, BrAC, generally reported in milligrams of alcohol per liter of breathed air. The relationship between BrAC and BAC is complex, and is affected by many factors.
Calibration
Calibration is the process of checking and adjusting the internal settings of a breath analyzer by comparing and adjting its test results to a known alcohol standard. Breath analyzer sensors drift over time and require periodic calibration to ensure accuracy. Many handheld breath analyzers sold to consumers use a silicon oxide sensor (also called a semiconductor sensor) to determine the alcohol concentration. These sensors are prone to contamination and interference from substances other than breath alcohol, and require recalibration or replacement every six months. Higher-end personal breath analyzers and professional-use breath alcohol testers use platinum fuel cell sensors. These too require recalibration but at less frequent intervals than semiconductor devices, usually once a year.
There are two ways of calibrating a precision fuel cell breath analyzer, the wet-bath and the dry-gas methods. Each method requires specialized equipment and factory-trained technicians. It is not a procedure that can be conducted by untrained users or without the proper equipment.
The dry-gas method utilizes a portable calibration standard which is a precise mixture of ethanol and inert nitrogen available in a pressurized canister. Initial equipment costs are less than alternative methods and the steps required are fewer. The equipment is also portable allowing calibrations to be done when and where required.
The wet-bath method utilizes an ethanol/water standard in a precise specialized alcohol concentration, contained and delivered in specialized breath simulator equipment. The wet-bath method has a higher initial cost and is not intended to be portable. The standard must be fresh and replaced regularly. In addition, the assumed water-air partition ratio for aqueous ethanol must be taken into account along with its associated uncertainty.
Some semiconductor models are designed specifically to allow the sensor module to be replaced without the need to send the unit to a calibration lab.
Non-specific analysis
One major problem with older breath analyzers is non-specificity: the machines identify not only the ethyl alcohol (or ethanol) found in alcoholic beverages but also other substances similar in molecular structure or reactivity, "interfering compounds".
The oldest breath analyzer models pass breath through a solution of potassium dichromate, which oxidizes ethanol into acetic acid, changing color in the process. A monochromatic light beam is passed through this sample, and a detector records the change in intensity and, hence, the change in color, which is used to calculate the percent alcohol in the breath. However, since potassium dichromate is a strong oxidizer, numerous alcohol groups can be oxidized by it, producing false positives. This source of false positives is unlikely as very few other substances found in exhaled air are oxidizable.
Infrared-based breath analyzers project an infrared beam of radiation through the captured breath in the sample chamber and detect the absorbance of the compound as a function of the wavelength of the beam, producing an absorbance spectrum that can be used to identify the compound, as the absorbance is due to the harmonic vibration and stretching of specific bonds in the molecule at specific wavelengths (see infrared spectroscopy). The characteristic bond of alcohols in infrared is the O-H bond, which gives a strong absorbance at a short wavelength. The more light is absorbed by compounds containing the alcohol group, the less reaches the detector on the other side—and the higher the reading. Other groups, most notably aromatic rings and carboxylic acids can give similar absorbance readings.
Some natural and volatile interfering compounds do exist, however. For example, the National Highway Traffic Safety Administration has found that dieters and diabetics may have acetone levels hundreds or even thousands of times higher than those in others. Acetone is one of the many substances that can be falsely identified as ethyl alcohol by some breath machines. However, fuel cell based systems are non-responsive to substances like acetone.
Substances in the environment can also lead to false BAC readings. For example, methyl tert-butyl ether, a common gasoline additive, has been alleged anecdotally to cause false positives in persons exposed to it. Tests have shown this to be true for older machines; however, newer machines detect this interference and compensate for it. Any number of other products found in the environment or workplace can also cause erroneous BAC results. These include compounds found in lacquer, paint remover, celluloid, gasoline, and cleaning fluids, especially ethers, alcohols, and other volatile compounds.
Pharmacokinetics
Absorption of alcohol continues for anywhere from 20 minutes (on an empty stomach) to two-and-one-half hours (on a full stomach) after the last consumption, generally taking around 40-50 minutes. During the absorptive phase, the concentration of alcohol throughout the body changes unpredictably, as it is affected by gastrointestinal physiology such as irregular contraction patterns. After absorption, the concentrations in the body settle down and follow predictable patterns. During absorption, the BAC in arterial blood will generally be higher than in venous blood, but post-absorption, venous BAC will be higher than arterial BAC. This is especially clear with bolus dosing, chugging a single large drink. With additional doses of alcohol, the definitions of absorption and post-absorption are less clear. However, once absorption of the last drink has finished, the concentrations will follow standard post-absorption curves. It is also not always clear from a BAC graph when the absorption phase finishes - for example, the body can reach a sustained equilibrium BAC where absorption and elimination are proportional.
Across all phases, BrAC correlates closely with arterial BAC. Arterial blood distributes oxygen throughout the body. Breath alcohol is a representation of the equilibrium of alcohol concentration as the blood gases (alcohol) pass from the arterial blood into the lungs to be expired in the breath. The ratio of ABAC:BrAC is 2294 ± 56 across all phases and 2251 ± 46 [2141-2307] in the post-absorption phase. For example, a breathalyzer measurement of 0.10 mg/L of breath alcohol characterises approximately 0.0001×2251 g/L, or 0.2251 g/L of arterial blood alcohol concentration (equivalent to 0.2251 permille or 0.02251% BAC).
The ratio of venous blood alcohol content to breath alcohol content may vary significantly, from 1300:1 to 3100:1. Assuming a blood-alcohol concentration of 0.07%, for example, a person could have a partition ratio of 1500:1 and a breath test reading of 0.10 g/2100 mL, over the legal limit in some jurisdictions. However, low partition ratios are generally observed during the absorption phase. Post-absorption, the ratio is relatively fixed, 2382 ± 119 [2125–2765], although this ratio was measured in a laboratory environment and variation may be larger in real-world scenarios.
Other false positives of high BrAC and also blood reading are related to patients with proteinuria and hematuria, due to kidney metabolization and failure. The metabolization rate of related patients with kidney damage is abnormal in relation to percent in alcohol in the breath. However, since potassium dichromate is a strong oxidizer, numerous alcohol groups can be oxidized by kidney and blood filtration, producing false positives.
Breathing pattern
It is sometimes said that the exhaled air analyzed by the breathalyzer is "alveolar air", coming from the alveoli in close proximity to the blood in pulmonary circulation and containing ethanol in concentrations proportional to that blood approximated by Henry's law. However, the alcohol in the exhaled air comes essentially from the airways of the lung, and not from the alveoli. The alcohol acts similarly to water vapor, so it is instructive to study the humidity of lung air. During breathing, the inspired air picks up water and alcohol from the airways. Almost all uptake occurs in the upper airways; thus, the BrAC is most affected by the alcohol concentration in the bronchial circulation, which supplies blood to these airways. When the air reaches the alveoli, it is already near equilibrium - this is why inhaling dry air does not dry out the lungs significantly. With exhalation, water and alcohol are rapidly lost to the airways, primarily within the fifth to fifteenth generations of branching. Nonetheless, as may be evidenced by seeing one's breath in the cold, some water vapor does not get re-absorbed by the airways and is exhaled, and similarly some alcohol is exhaled during breathing. But the relationship of the alcohol concentration of this air to the concentration of alcohol in the blood is somewhat suspect and can be affected by many variables.
As air is exhaled, the alcohol concentration of the exhaled air increases over time, rising significantly in the first few seconds and then slowing down after, but not leveling out until the subject stops exhaling. This is not because there is a "dead space" of non-alcoholic air in the airways - the alcohol concentration is nearly identical in all regions of the lung. Rather, it is because, during exhalation, water and alcohol are being redeposited on the airways, primarily the trachea and generations 6 though 12 of the airways. As more fluid is deposited on the mucous surfaces, the remaining fluid travels further, resulting in more alcohol being recorded by the breathalyzer. The recorded alcohol concentrations never reach the alveolar alcohol concentration, even if the subject exhales as deeply as possible. According to Henry's law, alveolar air alcohol concentration would be pulmonary BAC divided by 1756, compared to the BrAC which is arterial blood concentration divided by 2251. When the subject stops exhaling, the alcohol concentration levels off - this does not indicate that alveolar air has been obtained, as it will level off regardless of the point at which the subject stops exhaling. But it does mean that end-exhaled BrAC is readily obtained. This brings up the question of what is meant by reporting BrAC as a single number; is it the "deep-lung air", the highest possible reading obtainable by the subject's full exhalation? Or is it the zero concentration at the initial part of the curve? Hlastala suggests using the average BrAC during the exhalation, which corresponds to the BrAC measured at about the 5-second mark. The Supreme Court of California determined that the BrAC is defined as the alcohol concentration of the last part of the subject’s expired breath.
End-exhaled BrAC varies depending on several factors. Most alcohol breath testers require a minimum exhalation volume (normally between 1.1 and 1.5 L) or minimum six-second exhalation time before the breath sample is accepted. This raises concerns for subject with smaller lung volumes - they must exhale a greater fraction of their available lung volume compared with a larger subject. A mathematical model suggests that a 2L-lung-capacity subject's end-exhaled BrAC may read 35% higher than a 6L subject for the same minimum 1.5L exhalation and alveolar alcohol concentration. For exhalation to the maximum extent, such as under typical laboratory conditions, measured BrAC is unaffected by lung size. The subject's body temperature and breath temperature also influence results, with an increase in temperature corresponding to an increase in measured BrAC. Furthermore, the humidity and temperature of the ambient air can decrease results by as much as 10%. The result of these factors is that the breath test is more forgiving for some subjects than others. Nonetheless, the overall variance due to how much one breathes out is usually low, and some breathalyzers compensate for the volume of air.
Jones tested several breathing patterns immediately before and during breathalyzer use and found the following changes (in order of effect):
Hyperventilation by rapid inspiration and expiration of room air for 20 seconds before forced expiration - decrease by 10%
Moderate inspiration through mouth and deep expiration - control
Deep expiration without an inspiration - statistically insignificant increase
Inspiration through the nose before a deep expiration. - 1.3% increase
Deep inspiration followed by a slow (20 second) expiration. - 2.0% increase
Mouth closed for 5 minutes (shallow breathing) before nose-inspiration and a forced expiration. - 7.7% increase
Inspiration through the nose followed by breath-holding for 30 seconds before forced expiration. - 12.6% increase
A normal inspiration with breath-holding for 30 seconds before a forced expiration. - 15.7% increase
Overall, the results show an increase in measured BrAC with increased contact between the lungs and the measured air. Exercising immediately before the test, such as running up and down a flight of stairs, can also reduce measured BrAC by 13% or more, with the combined effect of exercise and hyperventilation reaching 20%.
Mouth alcohol
One of the most common causes of falsely high breath analyzer readings is the existence of mouth alcohol. In analyzing a subject's breath sample, the breath analyzer's internal computer is making the assumption that the alcohol in the breath sample came from the lungs. However, alcohol may have come from the mouth, throat or stomach for a number of reasons. A very tiny amount of alcohol from the mouth, throat or stomach can have a significant impact on the breath-alcohol reading.
Recent use of mouthwash or breath fresheners can also skew results upward, as they can contain fairly high levels of alcohol. Listerine mouthwash, for example, contains 26.9% alcohol, and can skew results for between 5 and 10 minutes. A scientist tested the effects of Binaca breath spray on an Intoxilyzer 5000. He performed 23 tests with subjects who sprayed their throats and obtained readings as high as 0.81—far beyond legal levels. The scientist also noted that the effects of the spray did not fall below detectable levels until after 18 minutes.
Other than those, the most common source of mouth alcohol is from belching or burping. This causes the liquids and/or gases from the stomach—including any alcohol—to rise up into the soft tissue of the esophagus and oral cavity, where it will stay until it has dissipated. The American Medical Association concludes in its Manual for Chemical Tests for Intoxication (1959): "True reactions with alcohol in expired breath from sources other than the alveolar air (eructation, regurgitation, vomiting) will, of course, vitiate the breath alcohol results." Acid reflux, or gastroesophageal reflux disease, can greatly exacerbate the mouth-alcohol problem. The stomach is normally separated from the throat by a valve, but when this valve becomes incompetent or herniated, there is nothing to stop the liquid contents in the stomach from rising and permeating the esophagus and mouth. The contents—including any alcohol—are then later exhaled into the breathalyzer. One study of 10 individuals suffering from this condition did not find any actual increase in breath ethanol.
Mouth alcohol can also be created in other ways. Dentures, some have theorized, will trap alcohol, although experiments have shown no difference if the normal 15 minute observation period is observed. Periodontal disease can also create pockets in the gums which will contain the alcohol for longer periods. Also known to produce false results due to residual alcohol in the mouth is passionate kissing with an intoxicated person.
To help guard against mouth-alcohol contamination, certified breath-test operators and police officers are trained to observe a test subject carefully for at least 15–20 minutes before administering the breath test. Some instruments also feature built-in safeguards. The Intoxilyzer 5000 features a "slope" parameter. This parameter detects any decrease in alcohol concentration of 0.006 g per 210 L of breath in 0.6 second, a condition indicative of residual mouth alcohol, and will result in an "invalid sample" warning to the operator, notifying the operator of the presence of the residual mouth alcohol. Other instruments require that the individual be tested twice at least two minutes apart. Mouthwash or other mouth alcohol will have somewhat dissipated after two minutes and cause the second reading to disagree with the first, requiring a retest. Many preliminary breath testers, however, feature no such safeguards.
Myths about accuracy
There are a number of substances or techniques that can supposedly "fool" a breath analyzer (i.e., generate a lower blood alcohol content).
A 2003 episode of the science television show MythBusters tested a number of methods that supposedly allow a person to fool a breath analyzer test. The methods tested included breath mints, onions, denture cream, mouthwash, pennies and batteries; all of these methods proved ineffective. The show noted that using these items to cover the smell of alcohol may fool a person, but, since they will not actually reduce a person's BrAC, there will be no effect on a breath analyzer test regardless of the quantity used, if any, it appeared that using mouthwash only raised the BrAC. Pennies supposedly produce a chemical reaction, while batteries supposedly create an electrical charge, yet neither of these methods affected the breath analyzer results.
The MythBusters episode also pointed out another complication: it would be necessary to insert the item into one's mouth (for example, eat an onion, rinse with mouthwash, conceal a battery), take the breath test, and then possibly remove the item — all of which would have to be accomplished discreetly enough to avoid alerting the police officers administering the test (who would obviously become very suspicious if they noticed that a person was inserting items into their mouth prior to taking a breath test). It would likely be very difficult, especially for someone in an intoxicated state, to be able to accomplish such a feat.
In addition, the show noted that breath tests are often verified with blood tests (BAC, which are more accurate) and that even if a person somehow managed to fool a breath test, a blood test would certainly confirm a person's guilt.
Other substances that might reduce the BrAC reading include a bag of activated charcoal concealed in the mouth (to absorb alcohol vapor), an oxidizing gas (such as N2O, Cl2, O3, etc.) that would fool a fuel cell type detector, or an organic interferent to fool an infrared absorption detector. The infrared absorption detector is more vulnerable to interference than a laboratory instrument measuring a continuous absorption spectrum since it only makes measurements at particular discrete wavelengths. However, due to the fact that any interference can only cause higher absorption, not lower, the estimated blood alcohol content will be overestimated. Additionally, Cl2 is toxic and corrosive.
A 2007 episode of the Spike network's show Manswers showed some of the more common and not-so-common ways of attempts to beat the breath analyzer, none of which work. Test 1 was to suck on a copper-coated coin such as a penny. Test 2 was to hold a battery on the tongue. Test 3 was to chew gum. None of these tests showed a "pass" reading if the subject had consumed alcohol.
Law enforcement
In general, two types of breathalyzer are used. Small hand-held breathalyzers are not reliable enough to provide evidence in court but reliable enough to justify an arrest. These devices may be used by officers in the field as a form of "field sobriety test" commonly called "preliminary breath test" or "preliminary alcohol screening", or as evidential devices in point of arrest testing. Larger breathalyzer devices found in police stations can be used to produce court evidence, These desktop analyzers generally use infrared spectrophotometer technology, electrochemical fuel cell technology, or a combination of the two.
All breath alcohol testers used by law enforcement in the United States of America must be approved by the Department of Transportation's National Highway Traffic Safety Administration.
Breath alcohol laws
The breath alcohol content reading may be used in prosecutions of the crime of driving under the influence of alcohol (sometimes referred to as driving or operating while intoxicated) in several ways. Historically, states in the US initially prohibited driving with a high level of BAC, and did not have any laws regarding BrAC. A BrAC test result was merely presented as indirect evidence of BAC. Where the defendant had refused to take a subsequent blood test, the only way the state could prove BAC was by presenting scientific evidence of how alcohol in the breath gets there from alcohol in the blood, along with evidence of how to convert from one to the other. DUI defense attorneys frequently contested the scientific reliability of such evidence. Before September 2011, South Dakota relied solely on blood tests to ensure accuracy.
States responded in different ways to the inability to rely on breathalyzer evidence. Many states such as California modified their statutes so to make a certain level of alcohol in the breath illegal per se. In other words, the BrAC level itself became the direct predicate evidence for conviction, with no need to estimate BAC. In per se jurisdictions such as the UK, it is automatically illegal to drive a vehicle with a sufficiently high breath alcohol concentration (BrAC). The breath analyzer reading of the operator will be offered as evidence of that crime, and challenges can only be offered on the basis of an inaccurate reading.
In other states, such as California and New Jersey, the statute remains tied to BAC, but the BrAC results of certain machines have been judicially deemed presumptively accurate substitutes for blood testing when used as directed. While BrAC tests are not necessary to prove a defendant was under the influence, laws in these states create a rebuttable presumption, which means it is presumed that the driver was intoxicated given a high BrAC reading, but that presumption can be rebutted if a jury finds it unreliable or if other evidence establishes a reasonable doubt as to whether the person actually drove with a breath or blood alcohol level of 0.08% or greater.
Another issue is that the BrAC is typically tested several hours after the time of driving. Some jurisdictions, such as the State of Washington, allow the use of breath analyzer test results without regard as to how much time passed between operation of the vehicle and the time the test was administered, or within a certain number of hours of testing. Other jurisdictions use retrograde extrapolation to estimate the BAC or BrAC at the time of driving.
One exception to criminal prosecution is the state of Wisconsin, where a first time drunk driving offense is normally a civil ordinance violation.
Breath levels
There is no international consensus on the statutory ratio of blood to breath levels, ranging from 2000:1 (most of Europe) to 2100:1 (US) to 2300:1 (UK). In the US, the ratio of 2100:1 was determined based on studies done in 1930-1950, with a 1952 report of the National Safety Council establishing the 2100:1 figure. The NSC has acknowledged that more recent research shows the actual relationship is most probably higher than 2100:1 and closer to 2300:1, but opines that this difference is of minimal practical significance in law enforcement. The use of the lower 2100:1 factor errs on the side of conservativism and can only favor the driver.
In early years, the range of the BrAC threshold in the US varied considerably between States. States have since adopted a uniform 0.08% BrAC level, due to federal guidelines. It is said that the federal government ensures the passage of the federal guidelines by tying traffic safety highway funds to compliance with federal guidelines on certain issues, such as the federal government ensuring that the legal drinking age be the age of 21 across the 50 states.
Police in Victoria, Australia, use breathalyzers that give a recognized 20% tolerance on readings. Noel Ashby, former Victoria Police Assistant Commissioner (Traffic & Transport), claims that this tolerance is to allow for different body types.
Preliminary breath tests
The preliminary breath test or preliminary alcohol screening test uses small hand-held breath analyzers (hand-held breathalyzers). (The terms "preliminary breath test" ("PBT") and "preliminary alcohol screening test" reference the same devices and functions.) They are generally based on electrochemical platinum fuel cell analysis. These units are similar to some evidentiary breathalyzers, but typically are not calibrated frequently enough for evidentiary purposes. The test device typically provides numerical blood alcohol content (BAC) readings, but its primary use is for screening. In some cases, the device even has "pass/fail" indicia. For example, in Canada, PST devices, called "alcohol screening devices" are set so that, from 0 to 49 mg% it shows digits, from 50 to 99 mg% it shows the word "warn" and 100 mg% and above it shows "fail". These preliminary breath tests are sometimes categorised as part of field sobriety testing, although it is not part of the series of performance tests generally associated with field sobriety tests (FSTs) or standard field sobriety tests (SFSTs).
In Canada, a preliminary non-evidentiary screening device can be approved by Parliament as an approved screening device. In order to demand a person produce a breathalyzer sample an officer must have "reasonable suspicion" that the person drove with more than 80 mg alcohol per 100 mL of blood. The demand must be within three hours of driving. Any driver that refuses can be charged under s.254 of the Criminal Code. With the legalization of cannabis, updates to the criminal code are proposed that will allow a breathalyzer test to be administered without suspicion of impairment.
The US National Highway Traffic Safety Administration maintains a Conforming Products List of breath alcohol devices approved for preliminary screening use. In the United States, the main use of the preliminary breath test (PBT) is to establish probable cause for arrest. All states have implied consent laws, which means that by applying for a driver's license, drivers are agreeing to take an evidentiary chemical test (blood, breath, or urine) after being arrested for a DUI. But in US law, the arrest and subsequent test may be invalidated if it is found that the arrest lacked probable cause. The PBT establishes a baseline alcohol level that the police officer may use to justify the arrest. The result of the PBT is not generally admissible in court, except to establish probable cause, although some states, such as Idaho, permit data or "readings" from hand-held preliminary breath testers or preliminary alcohol screeners to be presented as evidence in court. In states such as Florida and Colorado, there are no penalties for refusing a PBT. Police are not obliged to advise the suspect that participation in a FST, PBT, or other pre-arrest procedures is voluntary. In contrast, formal evidentiary tests given under implied consent requirements are considered mandatory.
Refusal to take a preliminary breath test in the State of Michigan subjects a non-commercial driver to a "civil infraction" fine, with no violation "points", but is not considered to be a refusal under the general "implied consent" law. In some states, the state may present evidence of refusal to take a field sobriety test in court, although this is of questionable probative value in a drunk driving prosecution.
Different requirements apply in many states to drivers under DUI probation, in which case participation in a preliminary breath test may be a condition of probation, and for commercial drivers under "drug screening" requirements. Some US states, notably California, have statutes on the books penalizing preliminary breath test refusal for drivers under 21; however the Constitutionality of those statutes has not been tested. (As a practical matter, most criminal lawyers advise suspects who refuse a preliminary breath test or preliminary alcohol screening to not engage in discussion or "justifying" the refusal with the police.)
Evidentiary breath tests
In Canada, an evidentiary breath instrument can be designated as an approved instrument. The US National Highway Traffic Safety Administration maintains a Conforming Products List of breath alcohol devices approved for evidentiary use, Infrared instruments are also known as "evidentiary breath testers" and generally produce court-admissible results.
Drinking after driving
A common defense to an impaired driving charge (in appropriate circumstances) is that the consumption of alcohol occurred subsequent to driving. The typical circumstance where this comes up is when a driver consumes alcohol after a road accident, as an affirmative defense. This closely relates to absorptive stage intoxication (or bolus drinking), except that the consumption of alcohol also occurred after driving. This defense can be overcome by retrograde extrapolation (infra), but complicates prosecution.
While jurisdictions that recognise absorptive stage intoxication as a defense would also accept a defense of consumption after driving, some jurisdictions penalise post-driving drinking. While laws regarding absorption of alcohol consumed before (or while) driving are generally per se, most statutes directed to post-driving consumption allow defenses for circumstances related to activity not related to. In Canada, it is illegal to be over the impaired driving limits within 3 hours of driving (given as 2 hours by CDN DOJ); however, the new law allows a "drinking after driving" defence in a situation where a driver had no reason to expect a demand by the police for breath testing. South Africa is more straightforward, with a separate penalty applied for consumption "After An Accident" until reported to the police and if so required, has been medically examined.
Retrograde extrapolation
The breath analyzer test is usually administered at a police station, commonly an hour or more after the arrest. Although this gives the BrAC at the time of the test, it does not by itself answer the question of what it was at the time of driving. The prosecution typically provides an estimated alcohol concentration at the time of driving utilizing retrograde extrapolation, presented by expert opinion. This involves projecting back in time to estimate the BrAC level at the time of driving, by applying the physiological properties of absorption and elimination rates in the human body.
Extrapolation is calculated using five factors and a general elimination rate of 0.015/hour.
Example Time of breath test-10:00pm...Result of breath test-0.080...Time of driving-9:00pm (stopped by officer)...Time of last drink-8:00pm...Last food-12:00pm. Using these facts, an expert can say the person's last drink was consumed on an empty stomach, which means absorption of the last drink (at 8:00) was complete within one hour-9:00. At the time of the stop, the driver is fully absorbed. The test result of 0.080 was at 10:00. So the one hour of elimination that has occurred since the stop is added in, making 0.080+0.015=0.095 the approximate breath alcohol concentration at the time of the stop.
Consumer use
Public breathalyzers are becoming a method for consumers to test themselves at the source of alcohol consumption. These are used in pubs, bars, restaurants, charities, weddings and all types of licensed events. As breathalyzer tests have increased risk of transmission of coronavirus, they were temporarily suspended from use in Sweden.
Breathalyzer sensors
Photovoltaic assay The photovoltaic assay, used only in the dated photoelectric intoximeter, is a form of breath testing rarely encountered today. The process works by using photocells to analyze the color change of a redox (oxidation-reduction) reaction. A breath sample is bubbled through an aqueous solution of sulfuric acid, potassium dichromate, and silver nitrate. The silver nitrate acts as a catalyst, allowing the alcohol to be oxidized at an appreciable rate. The requisite acidic condition needed for the reaction might also be provided by the sulfuric acid. In solution, ethanol reacts with the potassium dichromate, reducing the dichromate ion to the chromium (III) ion. This reduction results in a change of the solution's color from red-orange to green. The reacted solution is compared to a vial of non-reacted solution by a photocell, which creates an electric current proportional to the degree of the color change; this current moves the needle that indicates BAC. Like other methods, breath testing devices using chemical analysis are prone to false readings. Compounds that have compositions similar to ethanol, for example, could also act as reducing agents, creating the necessary color change to indicate increased BAC.
Infrared spectroscopy Infrared breathalyzers allow a high degree of specificity for ethanol. Typically evidential breath alcohol instruments in police stations will work on the principle of infrared spectroscopy.
Fuel cell Fuel cell gas sensors are based on the oxidation of ethanol to acetaldehyde on an electrode. The current produced is proportional to the amount of alcohol present. These sensors are very stable, typically requiring calibration every 6 months, and are the type of sensor usually found in roadside breath testing devices.
Semiconductor Semiconductor gas sensors are based on the increase in conductance of a tin oxide layer in the presence of a reducing gas such as vaporized ethanol. They are found in inexpensive breathalyzers and their stability is not as reliable as fuel cell instruments.
| Technology | Law enforcement equipment | null |
1592048 | https://en.wikipedia.org/wiki/Volcanic%20island | Volcanic island | Geologically, a volcanic island is an island of volcanic origin. The term high island can be used to distinguish such islands from low islands, which are formed from sedimentation or the uplifting of coral reefs (which have often formed on sunken volcanoes).
Definition and origin
There are a number of volcanic islands that rise no more than above sea level, often classified as islets or rocks, while some low islands, such as Banaba, Henderson Island, Makatea, Nauru, and Niue, rise over above sea level.
The two types of islands are often found in proximity to each other, especially among the islands of the South Pacific Ocean, where low islands are found on the fringing reefs that surround most volcanic islands. Volcanic islands normally rise above a hotspot or subduction zone.
Habitability
Volcanic islands usually range in size between . Islands above a certain size usually have fresh groundwater, while low islands often do not, so volcanic islands are more likely to be habitable.
Many volcanic islands emerge from the deep abyss of the ocean, and feature rough or mountainous landscapes in their interiors and a diverse array of summit elevations. Researchers have observed that the island will often be covered by dense tropical forest. These limit settlement on the interior of many islands, forcing communities to develop along the coast. Larger islands may have rivers, resulting in flood hazards. Rivers deliver sediment downstream, which can dominate the shape of the coast and contribute to erosion. Tall volcanic islands are often surrounded by protective fringing or barrier reefs, creating lagoons.
The unique geological and geographical characteristics of volcanic islands make them prone to many natural hazards, which are expected to worsen due to climate change. These include volcanic eruptions, earthquakes, tsunamis, landslides, and severe weather events like hurricanes or typhoons. Studies have highlighted the importance of implementing effective risk mitigation plans that include nature-based solutions to improve societal safety on these islands. These involve leveraging natural processes and ecosystems to reduce hazard impacts. This can include the restoration of natural barriers like mangroves or coral reefs that protect against tsunamis and storm surges or the maintenance of natural water catchments that can mitigate flood risks.
| Physical sciences | Oceanic and coastal landforms | Earth science |
1594239 | https://en.wikipedia.org/wiki/Addition%20principle | Addition principle | In combinatorics, the addition principle or rule of sum is a basic counting principle. Stated simply, it is the intuitive idea that if we have A number of ways of doing something and B number of ways of doing another thing and we can not do both at the same time, then there are ways to choose one of the actions. In mathematical terms, the addition principle states that, for disjoint sets A and B, we have , provided that the intersection of the sets is without any elements.
The rule of sum is a fact about set theory, as can be seen with the previously mentioned equation for the union of disjoint sets A and B being equal to |A| + |B|.
The addition principle can be extended to several sets. If are pairwise disjoint sets, then we have:This statement can be proven from the addition principle by induction on n.
Simple example
A person has decided to shop at one store today, either in the north part of town or the south part of town. If they visit the north part of town, they will shop at either a mall, a furniture store, or a jewelry store (3 ways). If they visit the south part of town then they will shop at either a clothing store or a shoe store (2 ways).
Thus there are possible shops the person could end up shopping at today.
Inclusion–exclusion principle
The inclusion–exclusion principle (also known as the sieve principle) can be thought of as a generalization of the rule of sum in that it too enumerates the number of elements in the union of some sets (but does not require the sets to be disjoint). It states that if A1, ..., An are finite sets, then
Subtraction principle
Similarly, for a given finite set S, and given another set A, if , then . To prove this, notice that by the addition principle.
Applications
The addition principle can be used to prove Pascal's rule combinatorially. To calculate , one can view it as the number of ways to choose k people from a room containing n children and 1 teacher. Then there are ways to choose people without choosing the teacher, and ways to choose people that includes the teacher. Thus .
The addition principle can also be used to prove the multiplication principle.
| Mathematics | Combinatorics | null |
1594286 | https://en.wikipedia.org/wiki/Rule%20of%20product | Rule of product | In combinatorics, the rule of product or multiplication principle is a basic counting principle (a.k.a. the fundamental principle of counting). Stated simply, it is the intuitive idea that if there are ways of doing something and ways of doing another thing, then there are ways of performing both actions.
Examples
In this example, the rule says: multiply 3 by 2, getting 6.
The sets {A, B, C} and {X, Y} in this example are disjoint sets, but that is not necessary. The number of ways to choose a member of {A, B, C}, and then to do so again, in effect choosing an ordered pair each of whose components are in {A, B, C}, is 3 × 3 = 9.
As another example, when you decide to order pizza, you must first choose the type of crust: thin or deep dish (2 choices). Next, you choose one topping: cheese, pepperoni, or sausage (3 choices).
Using the rule of product, you know that there are 2 × 3 = 6 possible combinations of ordering a pizza.
Applications
In set theory, this multiplication principle is often taken to be the definition of the product of cardinal numbers. We have
where is the Cartesian product operator. These sets need not be finite, nor is it necessary to have only finitely many factors in the product.
An extension of the rule of product considers there are different types of objects, say sweets, to be associated with objects, say people. How many different ways can the people receive their sweets?
Each person may receive any of the sweets available, and there are people, so there are ways to do this.
Related concepts
The rule of sum is another basic counting principle. Stated simply, it is the idea that if we have a ways of doing something and b ways of doing another thing and we can not do both at the same time, then there are a + b ways to choose one of the actions.
| Mathematics | Combinatorics | null |
1594759 | https://en.wikipedia.org/wiki/Leaf%20vegetable | Leaf vegetable | Leaf vegetables, also called leafy greens, pot herbs, vegetable greens, or simply greens, are plant leaves eaten as a vegetable, sometimes accompanied by tender petioles and shoots. Leaf vegetables eaten raw in a salad can be called salad greens.
Nearly one thousand species of plants with edible leaves are known. Leaf vegetables most often come from short-lived herbaceous plants, such as lettuce and spinach. Woody plants of various species also provide edible leaves.
The leaves of many fodder crops are also edible for humans, but are usually only eaten under famine conditions. Examples include alfalfa, clover, and most grasses, including wheat and barley. Food processing, such as drying and grinding into powder or pulping and pressing for juice, may involve these crop leaves in a diet.
Leaf vegetables contain many typical plant nutrients, but their vitamin K levels are particularly notable since they are photosynthetic tissues. Phylloquinone, the most common form of the vitamin, is directly involved in photosynthesis.
Nutrition
Spinach, as an example of a leaf vegetable, is low in calories and fat per calorie, and high in dietary fiber, vitamin C, pro-vitamin A carotenoids, folate, manganese and vitamin K.
The vitamin K content of leaf vegetables is particularly high since these are photosynthetic tissues, and phylloquinone is involved in photosynthesis. Accordingly, users of vitamin K antagonist medications, such as warfarin, must take special care to limit the consumption of leaf vegetables.
Preparation
If leaves are cooked for food, they may be referred to in the United States as boiled greens. Leaf vegetables may be stir-fried, stewed, steamed, or consumed raw. Leaf vegetables stewed with pork is a traditional dish in soul food and Southern U.S. cuisine. They are also commonly eaten in South Asian dishes such as saag. Leafy greens can be used to wrap other ingredients into an edible package like a tortilla. Many green leafy vegetables, such as lettuce or spinach, can also be eaten raw, for example, in sandwiches or salads. A green smoothie enables large quantities of raw leafy greens to be consumed by blending the leaves with fruit and water.
Africa
In certain countries of Africa, various species of nutritious amaranth are widely eaten boiled.
Celosia argentea var. argentea or "Lagos spinach" is one of the main boiled greens in West African cuisine.
Greece
In Greek cuisine, khorta (χόρτα, literally 'greens') are a typical side dish, eaten hot or cold and usually seasoned with olive oil and lemon.
At least 80 different kinds of greens are used, depending on the area and season, including black mustard, dandelion, wild sorrel, chicory, fennel, chard, kale, mallow, black nightshade, lamb's quarters, wild leeks, hoary mustard, charlock, smooth sow thistle and even the fresh leaves of the caper plant.
Italy
Preboggion, a mixture of different wild boiled leaf vegetables, is used in Ligurian cuisine to stuff ravioli and pansoti. One of the main ingredients of preboggion are borage (Borago officinalis) leaves.
Preboggion is also sometimes added to minestrone soup and frittata.
Poland
Botwinka (or boćwinka) is a soup that features beet stems and leaves as one of its main ingredients. The word "botwinka" is the diminutive form of "botwina" which refers to leafy vegetables like chard and beet leaves.
United States
In the cuisine of the Southern United States and traditional African-American cuisine, turnip, collard, kale, garden cress, dandelion, mustard, and pokeweed greens are commonly cooked and often served with pieces of ham or bacon. The boiling water, called potlikker, is used as broth. Water in which pokeweed has been prepared contains toxins that have been removed by boiling and should be discarded.
Sauteed escarole is a primary ingredient in the Italian-American dish Utica greens.
List of leaf vegetables
Agastache foeniculum — anise hyssop (western North America)
Allium fistulosum — Welsh onion (East Asia)
Alternanthera sissoo — sissoo spinach (Brazil)
Basella alba — Malabar spinach (India, Southeast Asia, New Guinea)
Beta vulgaris — beets, including beet greens, Swiss chard
Brassica oleracea — wild cabbage, including cabbage, gai lan, Jersey cabbage, kale, red cabbage, savoy cabbage, collard greens, mustard greens, kohlrabi and more
Brassica rapa — field mustard, including napa cabbage, bok choy, bomdong, choy sum, komatsuna, rapini, tatsoi, radish greens,and more
Campanula versicolor — various-colored bellflower (southeastern Italy to the Balkans)
Chenopodium quinoa — quinoa (western Andes of South America)
Cichorium endivia — endive, including escarole
Cichorium intybus — chicory (Europe)
Claytonia perfoliata — palsingat (western North America)
Cnidoscolus aconitifolius — chaya (Yucatán Peninsula of Mexico)
Daucus carota subsp. sativus — carrot (Europe and Southwestern Asia)
Eruca vesicaria — arugula or rocket (Mediterranean region)
Foeniculum vulgare — fennel (southern Europe)
Gynura bicolor — edible gynura (China, Thailand, Myanmar)
Gynura procumbens — longevity spinach (China, Southeast Asia, and Africa)
Hemerocallis fulva — orange day-lily (China or Japan)
Lepidium meyenii — maca (Andes)
Lactuca sativa — lettuce, including celtuce, iceberg lettuce, red leaf lettuce, romaine lettuce
Nasturtium officinale — watercress (Europe and Asia)
Malva moschata — musk mallow (Europe and southwestern Asia)
Moringa oleifera — moringa (Indian subcontinent)
Perilla frutescens — shisho perilla (Southeast Asia and Indian highlands)
Rumex acetosa — garden sorrel (most of Europe, temperate Asia, North America, and Greenland)
Sassafras albidum — sassafras (eastern North America)
Sauropus androgynus — katuk (South Asia and Southeast Asia)
Spinacia oleracea — spinach (central and western Asia)
Solanum aethiopicum — nakati (Asia and tropical Africa)
Trigonella foenum-graecum — fenugreek (India)
Tropaeolum majus — garden nasturtium (Andes)
Viola odorata — sweet violet (Europe, northern Africa, Syria)
Postharvest diseases
Postharvest diseases cause up to 50% losses of leaf vegetables. These are fungal, bacterial, and much less commonly viral. The most important remedy is temperature-controlled storage, although it is also important to prevent mechanical damage as this provides entryways for pathogens. Uncontaminated water for washing vegetables is of lesser but still significant importance.
Common bacterial pathogens include: Xanthomonas campestris pv. vitians, Pseudomonas viridiflava, P. cichorii, and P. marginalis, P. syringae pv. aptata, X. campestris pv. campestris, X. campestris pv. raphani, P. syringae pv. maculicola, P. syringae pv. alisalensis, Pectobacterium spp. including Pectobacterium carotovorum subsp. odoriferum and Pectobacterium aroidearum, Dickeya spp., Pseudomonas marginalis, and Pseudomonas viridiflava.
Common fungal pathogens include: Alternaria brassicicola, A. alternata, A. arborescens, A. tenuissima, A. japonica, Colletotrichum higginsianum, Colletotrichum dematium f. spinaciae, Microdochium panattonianum, Stemphylium botryosum, Cladosporium variabile, Cercospora beticola, C. lactucae-sativae, C. brassicicola, C. acetosella, Botrytis cinerea, Golovinomyces cichoracearum, Podosphaera fusca, Erysiphe cruciferarum, E. polygoni, E. heraclei, Sclerotinia sclerotiorum, and S. minor.
Common oomycete pathogens include: Albugo occidentalis, A. ipomoeae-aquaticae, A. candida, Hyaloperonospora parasitica, Bremia lactucae, Peronospora effusa, and Peronospora farinosa f.sp. betae.
Fungicides such as prochloraz can be used to manage some of these.
Gallery
| Biology and health sciences | Leafy vegetables | Plants |
1595197 | https://en.wikipedia.org/wiki/Sciadopitys | Sciadopitys | Sciadopitys, commonly called umbrella pines, is a genus of a unique conifers now endemic to Japan. The sole living member of the family Sciadopityaceae is Sciadopitys verticillata, a living fossil. The oldest fossils of Sciadopitys are from the Late Cretaceous of Japan, and the genus was widespread in Laurasia during most of the Cenozoic, especially in Europe until the Pliocene.
| Biology and health sciences | Pinophyta (Conifers) | Plants |
171865 | https://en.wikipedia.org/wiki/Agricultural%20policy | Agricultural policy | Agricultural policy describes a set of laws relating to domestic agriculture and imports of foreign agricultural products. Governments usually implement agricultural policies with the goal of achieving a specific outcome in the domestic agricultural product markets. Well designed agricultural policies use predetermined goals, objectives and pathways set by an individual or government for the purpose of achieving a specified outcome, for the benefit of the individual(s), society and the nations' economy at large. The goals could include issues such as biosecurity, food security, rural poverty reduction or increasing economic value through cash crop or improved food distribution or food processing.
Agricultural policies take into consideration the primary (production), secondary (such as food processing, and distribution) and tertiary processes (such as consumption and supply in agricultural products and supplies). Outcomes can involve, for example, a guaranteed supply level, price stability, product quality, product selection, land use or employment. Governments can use tools like rural development practices, agricultural extension, economic protections, agricultural subsidies or price controls to change the dynamics of agricultural production, or improve the consumer impacts of the production.
Agricultural policy has wide reaching primary and secondary effects. Agriculture has large impacts on climate change, estimated to be contributing 20–25% of global annual emissions as of 2010. Moreover, agricultural policy needs to account for a lot of shocks to the system: for example, agriculture is highly vulnerable to the negative impacts of climate change, such as decreases in water access, geophysical processes such as ocean level rise and changing weather, and socioeconomic processes that affect farmers, many of whom are in subsistence economic conditions. In order for global climate change mitigation and adaptation to be effective a wide range of policies need to be implemented to reduce the risk of negative climate change impacts on agriculture and greenhouse gas emissions from the agriculture sector.
Agriculture policy concerns
An example of the breadth and types of agriculture policy concerns can be found in the Australian Bureau of Agricultural and Resource Economics article "Agricultural Economies of Australia and New Zealand" which says that the major challenges and issues faced by their industrial agriculture industry are:
marketing challenges and consumer tastes
international trading environment (world market conditions, barriers to trade, quarantine and technical barriers, maintenance of global competitiveness and market image, and management of biosecurity issues affecting imports and the disease status of exports)
biosecurity (pests and diseases such as bovine spongiform encephalopathy (BSE), avian influenza, foot and mouth disease, citrus canker, and sugarcane smut)
infrastructure (such as transport, ports, telecommunications, energy and irrigation facilities)
management skills and labor supply (With increasing requirements for business planning, enhanced market awareness, the use of modern technology such as computers and global positioning systems and better agronomic management, modern farm managers will need to become increasingly skilled. Examples: training of skilled workers, the development of labor hire systems that provide continuity of work in industries with strong seasonal peaks, modern communication tools, investigating market opportunities, researching customer requirements, business planning including financial management, researching the latest farming techniques, risk management skills)
coordination (a more consistent national strategic agenda for agricultural research and development; more active involvement of research investors in collaboration with research providers developing programs of work; greater coordination of research activities across industries, research organisations and issues; and investment in human capital to ensure a skilled pool of research personnel in the future.)
technology (research, adoption, productivity, genetically modified (GM) crops, investments)
water (access rights, water trade, providing water for environmental outcomes, assignment of risk in response to the reallocation of water from consumptive to environmental use, accounting for the sourcing and allocation of water)
resource access issues (management of native vegetation, the protection and enhancement of biodiversity, sustainability of productive agricultural resources, and landholder responsibilities)
Poverty reduction
Policymakers working in poverty reduction in the agriculture sector assess, plan, or enact policies aimed to address the needs of persons living in poverty. Agriculture has been a critical driver of poverty reduction in most developing countries, particularly in rural areas. Approximately 80% of the world's impoverished population, who primarily reside in rural areas and earn their livelihood through farming, can benefit from agriculture in terms of poverty reduction, income generation, and food security. Fostering agricultural development is therefore a crucial element of agricultural policy in a developing country. In addition, a recent Natural Resource Perspective paper by the Overseas Development Institute found that good infrastructure, education and effective information services in rural areas were necessary to improve the chances of making agriculture work for the poor.
During the 1980s and 1990s, there was a disregard for the agriculture sector among policymakers and investors, only regaining interest when the prices of staple food crops experienced a significant increase in the mid-2000s. As a result of agricultural policy neglect, there has been a scarcity of investment in infrastructure, which has hindered agricultural development and public goods, such as education, research and development and technology. Rural productive sectors and small agricultural enterprises suffer from market failures due to policies favouring urban areas and lending policies biased against small-scale agricultural firms. Neglect in implementing agriculture policy has been detected in several developing countries. In Indonesia, since the Asian Financial Crisis of 1997 to 1998, the government's agricultural policy has been closely concentrated on achieving price stability and self-sufficiency for import-competing commodities, such as palm oil, sugar and rice.
International agencies such as the Food and Agriculture Organization (FAO), the World Bank, the International Fund for Agricultural Development (IFAD) and the Organisation for Economic Co-operation and Development (OECD), espouse the prioritisation of agricultural endeavours to support poverty reduction. The impact of agricultural policy on reducing poverty differs across countries and is influenced by a variety of factors, such as the level of government policy support, the degree of public and private investment in agriculture, the different types of agriculture, and the growth rates of agriculture parallel to non-agriculture sectors. In particular, investment in agricultural research and development has been shown to be highly influential on agricultural GDP growth and poverty reduction. Government policies play a key role in promoting agricultural activities, such as irrigation systems, roads and telecommunication systems, land reform, power in rural areas, fiscal support for research and development, pricing policies, assistance for new technologies, and markets for agricultural produce. Agricultural policies have contributed to meeting the goals related to increasing, diversifying, and improving agricultural production.
Agricultural policies aimed at reducing poverty include India's Pradhan Mantri Fasal Bima Yojana, which offers crop insurance to farmers to protect them from weather-related uncertainties and potential crop failures. This initiative provides farmers with financial aid for crop loss, reducing the risk of falling into poverty. Rwanda's Crop Intensification Program is another example of such policy, which provides farmers with inputs like fertilisers, improved seeds, and pesticides, as well as training and technical support to help them adopt more efficient farming practices. However, for agricultural policies to contribute to poverty reduction, it is essential that they collaborate effectively and cohesively with other sectors, such as tourism, sustainable economy, and industry.
Biosecurity
The biosecurity concerns facing industrial agriculture can be illustrated by:
the threat to poultry and humans from H5N1; possibly caused by the use of animal vaccines
the threat to cattle and humans from bovine spongiform encephalopathy (BSE); possibly caused by the unnatural feeding of cattle to cattle to minimize costs
the threat to industry profits from diseases like foot-and-mouth disease and citrus canker which increasing globalization makes harder to contain
Avian influenza
The use of animal vaccines can create new viruses that kill people and cause flu pandemic threats. H5N1 is an example of where this might have already occurred. According to the CDC article "H5N1 Outbreaks and Enzootic Influenza" by Robert G. Webster et al.: "Transmission of highly pathogenic H5N1 from domestic poultry back to migratory waterfowl in western China has increased the geographic spread. The spread of H5N1 and its likely reintroduction to domestic poultry increase the need for good agricultural vaccines. In fact, the root cause of the continuing H5N1 pandemic threat may be the way the pathogenicity of H5N1 viruses is masked by co-circulating influenza viruses or bad agricultural vaccines." Robert Webster explains: "If you use a good vaccine you can prevent the transmission within poultry and to humans. But if they have been using vaccines now [in China] for several years, why is there so much bird flu? There is bad vaccine that stops the disease in the bird but the bird goes on pooping out the virus and maintaining it and changing it. And I think this is what is going on in China. It has to be. Either there is not enough vaccine being used or there is substandard vaccine being used. Probably both. It's not just China. We can't blame China for substandard vaccines. I think there are substandard vaccines for influenza in poultry all over the world."
In response to the same concerns, Reuters reports Hong Kong infectious disease expert Lo Wing-lok indicating that vaccines have to take top priority. Julie Hall, who is in charge of the WHO's outbreak response in China, claimed that China's vaccinations might be masking the virus. The BBC reported that Wendy Barclay, a virologist at the University of Reading, UK said: "The Chinese have made a vaccine based on reverse genetics made with H5N1 antigens, and they have been using it. There has been a lot of criticism of what they have done because they have protected their chickens against death from this virus but the chickens still get infected, and then you get the drift - the virus mutates in response to the antibodies - and now we have a situation where we have five or six 'flavours' of H5N1 out there."
Bovine spongiform encephalopathy
Bovine spongiform encephalopathy (BSE), commonly known as "mad cow disease", is a fatal, neurodegenerative disease of cattle, which infects by a mechanism that surprised biologists upon its discovery in the late 20th century. In the UK, the country worst affected, 179,000 cattle were infected and 4.4 million were killed as a precaution. The disease can be transmitted to human beings who eat or inhale material from infected carcasses. In humans, it is known as new variant Creutzfeldt–Jakob disease (vCJD or nvCJD), and by June 2007, it had killed 165 people in Britain, and six elsewhere with the number expected to rise because of the disease's long incubation period. Between 460,000 and 482,000 BSE-infected animals had entered the human food chain before controls on high-risk offal were introduced in 1989.
A British inquiry into BSE concluded that the epidemic was caused by feeding cattle, who are normally herbivores, the remains of other cattle in the form of meat and bone meal (MBM), which caused the infectious agent to spread. The origin of the disease itself remains unknown. The current scientific view is that infectious proteins called prions developed through spontaneous mutation, probably in the 1970s, and there is a possibility that the use of organophosphorus pesticides increased the susceptibility of cattle to the disease. The infectious agent is distinctive for the high temperatures it is able to survive; this contributed to the spread of the disease in Britain, which had reduced the temperatures used during its rendering process. Another contributory factor was the feeding of infected protein supplements to very young calves instead of milk from their mothers.
Foot-and-mouth disease
Foot-and-mouth disease is a highly contagious and sometimes fatal viral disease of cattle and pigs. It can also infect deer, goats, sheep, and other bovids with cloven hooves, as well as elephants, rats, and hedgehogs. Humans are affected only very rarely. FMD occurs throughout much of the world, and while some countries have been free of FMD for some time, its wide host range and rapid spread represent cause for international concern. In 1996, endemic areas included Asia, Africa, and parts of South America. North America, Australia, New Zealand and Japan have been free of FMD for many years. Most European countries have been recognized as free, and countries belonging to the European Union have stopped FMD vaccination.
Infection with foot-and-mouth disease tends to occur locally, that is, the virus is passed on to susceptible animals through direct contact with infected animals or with contaminated pens or vehicles used to transport livestock. The clothes and skin of animal handlers such as farmers, standing water, and uncooked food scraps and feed supplements containing infected animal products can harbor the virus as well. Cows can also catch FMD from the semen of infected bulls. Control measures include quarantine and destruction of infected livestock, and export bans for meat and other animal products to countries not infected with the disease.
Because FMD rarely infects humans but spreads rapidly among animals, it is a much greater threat to the agriculture industry than to human health. Farmers around the world can lose huge amounts of money during a foot-and-mouth epidemic, when large numbers of animals are destroyed and revenues from milk and meat production go down. One of the difficulties in vaccinating against FMD is the huge variation between and even within serotypes. There is no cross-protection between serotypes (meaning that a vaccine for one serotype won't protect against any others) and in addition, two strains within a given serotype may have nucleotide sequences that differ by as much as 30% for a given gene. This means that FMD vaccines must be highly specific to the strain involved. Vaccination only provides temporary immunity that lasts from months to years. Therefore, rich countries maintain a policy of banning imports from all countries, not proven FMD-free by US or EU standards. This is a point of contention.
Although this disease is not dangerous to humans and rarely fatal to otherwise healthy animals, it reduces milk and meat production. Outbreaks can be stopped quickly if farmers and transporters are forced to abide by existing rules. Therefore, (besides temporary discomfort to the animals), any outbreak in the rich world should not be much more as a localized, cyclical economic problem. For countries with free roaming wildlife it is nearly impossible to prove that they are entirely free of this disease. If they try they are forced to erect nationwide fences, which destroys wildlife migration. Because detecting and reporting of FMD have enormously improved and sped up, almost all poor countries could now safely create FMD-free export zones. But rich countries refuse to change the rules. In effect, many poor tropical countries have no chance to meet current rules, so they are still today banned from exporting meat, even if many of them are FMD-free.
The result is that if drought hits, the poor try to cope by selling their few animals. This quickly saturates regional demand. The export ban then destroys the value of these animals, in effect destroying the most important coping mechanism of several hundreds of millions extremely poor households. The rules around meat exports have been changed many times, always to accommodate changing circumstances in rich countries, usually further reducing meat export chances for poor countries. For that reason, Kanya and many other countries find the rules very unjust. They are however discouraged to file a formal complaint with WTO by diplomats from rich countries.
Citrus canker
Citrus canker is a disease affecting citrus species that is caused by the bacterium Xanthomonas axonopodis. The infection causes lesions on the leaves, stems, and fruit of citrus trees, including lime, oranges, and grapefruit. While not harmful to humans, canker significantly affects the vitality of citrus trees, causing leaves and fruit to drop prematurely; a fruit infected with canker is safe to eat but too unsightly to be sold. The disease, which is believed to have originated in South East Asia, is extremely persistent when it becomes established in an area, making it necessary for all citrus orchards to be destroyed for the successful eradication of the disease. Australia, Brazil and the United States are currently experiencing canker outbreaks.
The disease can be detected in orchards and on fruit by the appearance of lesions. Early detection is critical in quarantine situations. Bacteria are tested for pathogenicity by inoculating multiple citrus species with the bacterium. Simultaneously, other diagnostic tests (antibody detection, fatty-acid profiling, and genetic procedures using PCR) are conducted to identify the particular canker strain. Citrus canker outbreaks are prevented and managed in a number of ways. In countries that do not have canker, the disease is prevented from entering the country by quarantine measures. In countries with new outbreaks, eradication programs that are started soon after the disease has been discovered have been successful; such programs rely on the destruction of affected orchards. When eradication has been unsuccessful and the disease has become established, management options include replacing susceptible citrus cultivars with resistant cultivars, applying preventive sprays of copper-based bactericides, and destroying infected trees and all surrounding trees within an appropriate radius.
The citrus industry is the largest fresh-fruit exporting industry in Australia. Australia has had three outbreaks of citrus canker; all three were successfully eradicated. The disease was found twice during the 1900s in the Northern Territory and was eradicated each time. During the first outbreak in 1912, every citrus tree north of latitude 19° south was destroyed, taking 11 years to eradicate the disease. In 2004, Asiatic citrus canker was detected in an orchard in Emerald, Queensland, and was thought to have occurred from the illegal import of infected citrus plants. The state and federal governments have ordered that all commercial orchards, all non-commercial citrus trees, and all native lime trees (C. glauca) in the vicinity of Emerald be destroyed rather than trying to isolate infected trees.
Food security
The United Nations Food and Agriculture Organization (FAO) defines food security as existing when "all people, at all times, have physical and economic access to sufficient safe and nutritious food that meets their dietary needs and food preferences for an active and healthy life". The four qualifications that must be met for a food secure system include physical availability, economic and physical access, appropriate utilization, and stability of the prior three elements over time.
Of the 6.7 billion people on the planet, about 2 billion are food insecure. As the global population grows to 9 billion by 2050, and as diets shift to emphasize higher energy products and greater overall consumption, food systems will be subjected to even greater pressure. Climate change presents additional threats to food security, affecting crop yields, distribution of pests and diseases, weather patterns, and growing seasons around the world.
Food security has thus become an increasingly important topic in agricultural policy as decision makers attempt to reduce poverty and malnutrition while augmenting adaptive capacity to climate change. The Commission on Sustainable Agriculture and Climate Change listed high-priority policy actions to address food security, including integrating food security and sustainable agriculture into global and national policies, significantly raising the level of global investment in food systems, and developing specific programs and policies to support the most vulnerable populations (namely, those that are already subject to food insecurity).
Food sovereignty
'Food sovereignty', a term coined by members of Via Campesina in 1996, is about the right of peoples to define their own food systems. Advocates of food sovereignty put the people who produce, distribute, and consume food at the centre of decisions on food systems and policies, rather than the demands of markets and corporations that they believe have come to dominate the global food system. This movement is advocated by a number of farmers, peasants, pastoralists, fisherfolk, indigenous peoples, women, rural youth, and environmental organizations.
Policy tools
An agricultural subsidy is a governmental subsidy paid to farmers and agribusinesses to manage the agricultural industry as one part of the various methods a government uses in a mixed economy. The conditions for payment and the reasons for the individual specific subsidies vary with farm product, size of the farm, nature of ownership, and country among other factors. Enriching peanut farmers for political purposes, keeping the price of a staple low to keep the poor from rebelling, stabilizing the production of a crop to avoid famine years, encouraging diversification and many other purposes have been suggested as the reason for specific subsidies.
Price floors or price ceilings set a minimum or maximum price for a product. Price controls encourage more production by a price floor or less production by a price ceiling.
A government can erect trade barriers to limit the number of goods imported (in the case of a Quota Share) or enact tariffs to raise the domestic price of imported products. These barriers give preference to domestic producers.
Objectives of market intervention
National security
Some argue that nations have an interest in assuring there is sufficient domestic production capability to meet domestic needs in the event of a global supply disruption. Significant dependence on foreign food producers makes a country strategically vulnerable in the event of war, blockade or embargo. Maintaining adequate domestic capability allows for food self-sufficiency that lessens the risk of supply shocks due to geopolitical events. Agricultural policies may be used to support domestic producers as they gain domestic and international market share. This may be a short term way of encouraging an industry until it is large enough to thrive without aid. Or it may be an ongoing subsidy designed to allow a product to compete with or undercut the foreign competition. This may produce a net gain for a government despite the cost of interventions because it allows a country to build up an export industry or reduce imports. It also helps to form the nation's supply and demand market.
Environmental protection and land management
Farm or undeveloped land composes the majority of land in most countries. Policies may encourage some land uses rather than others in the interest of protecting the environment. For instance, subsidies may be given for particular farming methods, forestation, land clearance, or pollution abatement.
Rural poverty and poverty relief
Subsidising farming may encourage people to remain on the land and obtain some income. This might be relevant to an agrarian country with many peasant farmers, but it may also be a consideration to more developed countries such as Poland. It has a very high unemployment rate, much farmland and retains a large rural population growing food for their own use.
Price controls may also be used to assist poor citizens. Many countries have used this method of welfare support as it delivers cheap food to the poorest in urban areas without the need to assess people to give them financial aid. This often goes at the cost of the rural poor, who then earn less from what is often their only realistic or potential source of income: agriculture. Because in almost all countries the rural poor are poorer than the urban poor, cheap food policies through price controls often increase overall poverty.
The same often counts for poverty relief in the form of food aid, which (unless while during severe drought) drives small producers in poor countries out of production. It tends to benefit lower middle class groups (sub-urban and urban poor) at the expense of the poorest 20 percent, who as a result remain deprived of customers.
Organic farming assistance
Welfare economics theory holds that sometimes private activities can impose social costs upon others. Industrial agriculture is widely considered to impose social costs through pesticide pollution and nitrate pollution. Further, agriculture uses large amounts of water, a scarce resource. Some economists argue that taxes should be levied on agriculture, or that organic agriculture, which uses little pesticides and experiences relatively little nitrate runoff, should be encouraged with subsidies. In the United States, 65% of the approximately $16.5 billion in annual subsidies went to the top 10% of farmers in 2002 because subsidies are linked to certain commodities. On the other hand, organic farming received $5 million for help in certification and $15 million for research over a 5-year time period.
Fair trade
Some advocate Fair Trade rules to ensure that poor farmers in developing nations that produce crops primarily for export are not exploited or negatively impacted by trade policies, practices, tariffs, and agreements which benefit one competitor at the expense of another - which advocates consider a dangerous "race to the bottom" in agricultural labor and safety standards. Opponents point out that most agriculture in developed nations is produced by industrial corporations (agribusiness) which are hardly deserving of sympathy, and that the alternative to exploitation is poverty.
Fair trade steak? Much of what developing countries export to the rich world, also comes from industrial corporations. The reason for that is, that rich countries have put up elaborate quality demands, most of whom make no factual health contribution. Small farmers often in effect meet these demands, but are rarely able to prove that in western standards. Therefore, the biggest impediment to the growth of small farming and therefore of fair trade in sectors beyond coffee and bananas, is these quality demands from the rich world.
Arguments against market intervention
Dumping of agricultural surpluses
In international trade parlance, when a company from country A sells a commodity below the cost of production into country B, this is called "dumping". A number of countries that are signatories to multilateral trade agreements have provisions that prohibit this practice.
When rich countries subsidize domestic production, the excess output is often given to the developing world as foreign aid. This process eliminates the domestic market for agricultural products in the developing world because the products can be obtained for free from western aid agencies. In developing nations where these effects are most severe, small farmers could no longer afford basic inputs and were forced to sell their land.
"Consider a farmer in Ghana who used to be able to make a living growing rice. Several years ago, Ghana was able to feed and export their surplus. Now, it imports rice. From where? Developed countries. Why? Because it's cheaper. Even if it costs the rice producer in the developed world much more to produce the rice, he doesn't have to make a profit from his crop. The government pays him to grow it, so he can sell it more cheaply to Ghana than the farmer in Ghana can. And that farmer in Ghana? He can't feed his family anymore." (Lyle Vanclief, former Canadian Minister of Agriculture [1997-2003])
According to the Institute for Agriculture and Trade Policy, corn, soybeans, cotton, wheat and rice are sold below the cost of production, or dumped. Dumping rates are approximately forty percent for wheat, between twenty-five and thirty percent for corn (maize), approximately thirty percent for soybeans, fifty-seven percent for cotton, and approximately twenty percent for rice. For example, wheat is sold for forty percent below cost.
According to Oxfam, "If developed nations eliminated subsidy programs, the export value of agriculture in lesser developed nations would increase by 24 %, plus a further 5.5 % from tariff equilibrium. ... exporters can offer US surpluses for sale at prices around half the cost of production; destroying local agriculture and creating a captive market in the process." Free trade advocates desire the elimination of all market distorting mechanisms (subsidies, tariffs, regulations) and argue that, as with free trade in all areas, this will result in aggregate benefit for all. This position is particularly popular in competitive agricultural exporting nations in both the developed and developing world, some of whom have banded together in the Cairns Group lobby. Canada's Department of Agriculture estimates that developing nations would benefit by about $4 billion annually if subsidies in the developed world were halved.
Agricultural independence
Many developing countries do not grow enough food to feed their own populations. These nations must buy food from other countries. Lower prices and free food save the lives of millions of starving people, despite the drop in food sales of the local farmers. A developing nation could use new improved farming methods to grow more food, with the ultimate goal of feeding their nation without outside help. New greenhouse methods, hydroponics, fertilizers, R/O water processors, hybrid crops, fast-growing hybrid trees for quick shade, interior temperature control, greenhouse or tent insulation, autonomous building gardens, sun lamps, mylar, fans, and other cheap tech can be used to grow crops on previously unarable land, such as rocky, mountainous, desert, and even Arctic lands. More food can be grown, reducing dependency on other countries for food.
Replacement crops can also make nations agriculturally independent. Sugar, for example, comes from sugar cane imported from Polynesia. Instead of buying the sugar from Polynesia, a nation can make sugar from sugar beets, maple sap, or sweetener from stevia plant, keeping the profits circulating within the nation's economy. Paper and clothes can be made of hemp instead of trees and cotton. Tropical foods won't grow in many places in Europe, but they will grow in insulated greenhouses or tents in Europe. Soybean plant cellulose can replace plastic (made from oil). Ethanol from farm waste or hempseed oil can replace gasoline. Rainforest medicine plants grown locally can replace many imported medicines. Alternates of cash crops, like sugar and oil replacements, can reduce farmers' dependency on subsidies in both developed and developing nations.
Market interventions may increase the cost to consumers for agricultural products, either via hidden wealth-transfers via the government, or increased prices at the consumer level, such as for sugar and peanuts in the US. This has led to market distortions, such as food processors using high fructose corn syrup as a replacement for sugar. High fructose corn syrup may be an unhealthy food additive, and, were sugar prices not inflated by government fiat, sugar might be preferred over high fructose corn syrup in the marketplace.
Developed world cases
Agriculture policy design strategy and examples
The concerns of agricultural policies are extensive, and includes ensuring the hygiene of salads, globalization management, and other emerging issues. The majority of the concerns fall into three categories: food supply for a growing population, livelihood insurance for farmers, and environmental protection. The theme of all approaches aiming to address these 3 types of concerns is to have a holistic view of their effects and externalities (a by-product of an action that affects others without their consent), because some policies intended to address one aspect of the concerns may have unintended harmful consequences that worsen other aspects while some have zero or negative beneficial effects. For example, subsidizing agricultural companies allows them to expand their industry and offer their products at lower prices to customers, but increases the firm's water and land usage which are at the cost of natural habitats. From an opposite perspective, if we protect the natural habitats and tax the agricultural firm for turning natural lands into factories, the prices of their products increase, making the firm's products too expensive for some customers. These externalities and trade-offs put the policymakers in a dilemma because our current global agriculture system is vulnerable to many disruptions such as weather changes, locality, manpower shifts, etc. Consequently, before we resolve this primary fragility of our agriculture system. It's of high cruciality for policymakers to weigh the trade-offs and adopt the most appropriate policies.
There are examples of the agricultural policy design mentioned above that are made by worldwide unions, countries, and states. While every specific situation requires its own specific agricultural policy design, these examples can provide useful models, insights, and lessons for future policymakers' reference and inspiration. The Common Agricultural Policy, published by E.U., uses government subsidies to encourage food production and farming industrialization in its early stage. In some areas, food production boomed so much that enormous food waste became a new problem. With food waste, the market was thrown into imbalance. Consequently, the price drop cost the farmers' utility and has led to a future reform known as the Marsholt Plan. Marsholt Plan and following reforms generally adjusted the agriculture market back to balance. Later reforms managed to spread the fund to farmers and increase each individual farmers' welfare instead of merely expanding croplands and industries. Starting from 2003, the Common Agricultural Policy fund is further detailed into individuals and environmental protection is finally put into consideration.
Overview: Europe and America
The farmer population is approximately five percent of the total population in the E.U. and 1.7% in the U.S. The total value of agricultural production in the E.U. amounted to 128 billion euros (1998). About forty-nine percent of this amount was accounted for by political measures: 37 billion euros due to direct payments and 43 billion euros from consumers due to the artificially high price. Eighty percent of European farmers receive a direct payment of 5,000 euros or less, while 2.2% receive a direct payment above 50,000 euros, totaling forty percent of all direct subsidies.
The average U.S. farmer receives $16,000 in annual subsidies. Two-thirds of farmers receive no direct payments. Of those that do, the average amount amongst the lowest paid eighty percent was $7000 from 1995 to 2003. Subsidies are a mix of tax reductions, direct cash payments and below-market prices on water and other inputs. Some claim that these aggregate figures are misleading because large agribusinesses, rather than individual farmers, receive a significant share of total subsidy spending. The Federal Agriculture Improvement and Reform Act of 1996 reduced farm subsidies, providing fixed payments over a period and replacing price supports and subsidies. The Farm Security and Rural Investment Act of 2002 contains direct and countercyclical payments designed to limit the effects of low prices and yields.
In the EU, €54 billion of subsidies are paid every year. An increasing share of the subsidies is being decoupled from production and lumped into the Single Farm Payment. While this has diminished the distortions created by the Common Agricultural Policy, many critics argue that a greater focus on the provision of public goods, such as biodiversity and clean water, is needed. The next major reform is expected for 2014 when a new long-term EU budget is coming into effect.
Environmental programs
The U.S. Conservation Reserve Program leases land from producers that take marginal land out of production and convert it back to a near-natural state by planting native grasses and other plants. The U.S. Environmental Quality Incentives Program subsidizes improvements which promote water conservation and other measures. This program is conducted under a bidding process using a formula where farmers request a certain percentage of cost share for improvement such as drip irrigation. Producers that offer the most environmental improvement based on a point system for the least cost are funded first. The process continues until that year's allocated funds are expended.
World Trade Organization actions
In April 2004 the World Trade Organization (WTO) ruled that 3 billion dollars in US cotton subsidies violate trade agreements and that almost 50% of EU sugar exports are illegal. In 1997–2003, US cotton exports were subsidized by an average of 48%. The WTO has extracted commitments from the Philippines government, making it lower import barriers to half their present levels over a span of six years, and allowing in drastically increased competition from the industrialised and heavily subsidised farming systems of North America and Europe. A recent Oxfam report estimated that average household incomes of maize farmers will be reduced by as much as 30% over the six years as cheap imports from the US drive down prices in the local markets. The report estimates that in the absence of trade restrictions, US subsidised maize could be marketed at less than half the price of maize grown on the Philippine island of Mindanao; and that the livelihoods of up to half a million Filipino maize farmers (out of the total 1.2 million) are under immediate threat.
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171879 | https://en.wikipedia.org/wiki/Brachistochrone%20curve | Brachistochrone curve | In physics and mathematics, a brachistochrone curve (), or curve of fastest descent, is the one lying on the plane between a point A and a lower point B, where B is not directly below A, on which a bead slides frictionlessly under the influence of a uniform gravitational field to a given end point in the shortest time. The problem was posed by Johann Bernoulli in 1696.
The brachistochrone curve is the same shape as the tautochrone curve; both are cycloids. However, the portion of the cycloid used for each of the two varies. More specifically, the brachistochrone can use up to a complete rotation of the cycloid (at the limit when A and B are at the same level), but always starts at a cusp. In contrast, the tautochrone problem can use only up to the first half rotation, and always ends at the horizontal. The problem can be solved using tools from the calculus of variations and optimal control.
The curve is independent of both the mass of the test body and the local strength of gravity. Only a parameter is chosen so that the curve fits the starting point A and the ending point B. If the body is given an initial velocity at A, or if friction is taken into account, then the curve that minimizes time differs from the tautochrone curve.
History
Galileo's problem
Earlier, in 1638, Galileo Galilei had tried to solve a similar problem for the path of the fastest descent from a point to a wall in his Two New Sciences. He draws the conclusion that the arc of a circle is faster than any number of its chords,From the preceding it is possible to infer that the quickest path of all [lationem omnium velocissimam], from one point to another, is not the shortest path, namely, a straight line, but the arc of a circle.
...
Consequently the nearer the inscribed polygon approaches a circle the shorter the time required for descent from A to C. What has been proven for the quadrant holds true also for smaller arcs; the reasoning is the same.
Just after Theorem 6 of Two New Sciences, Galileo warns of possible fallacies and the need for a "higher science". In this dialogue Galileo reviews his own work. Galileo studied the cycloid and gave it its name, but the connection between it and his problem had to wait for advances in mathematics.
Galileo’s conjecture is that “The shortest time of all [for a movable body] will be that of its fall along the arc ADB [of a quarter circle] and similar properties are to be understood as holding for all lesser arcs taken upward from the lowest limit B.”
In Fig.1, from the “Dialogue Concerning the Two Chief World Systems”, Galileo claims that the body sliding along the circular arc of a quarter circle, from A to B will reach B in less time than if it took any other path from A to B. Similarly, in Fig. 2, from any point D on the arc AB, he claims that the time along the lesser arc DB will be less than for any other path from D to B. In fact, the quickest path from A to B or from D to B, the brachistochrone, is a cycloidal arc, which is shown in Fig. 3 for the path from A to B, and Fig.4 for the path from D to B, superposed on the respective circular arc.
Introduction of the problem
Johann Bernoulli posed the problem of the brachistochrone to the readers of Acta Eruditorum in June, 1696. He said:
Bernoulli wrote the problem statement as:
{{Quote
|text=Given two points A and B in a vertical plane, what is the curve traced out by a point acted on only by gravity, which starts at A and reaches B in the shortest time.}}
Johann and his brother Jakob Bernoulli derived the same solution, but Johann's derivation was incorrect, and he tried to pass off Jakob's solution as his own. Johann published the solution in the journal in May of the following year, and noted that the solution is the same curve as Huygens' tautochrone curve. After deriving the differential equation for the curve by the method given below, he went on to show that it does yield a cycloid. However, his proof is marred by his use of a single constant instead of the three constants, vm, 2g and D, below.
Bernoulli allowed six months for the solutions but none were received during this period. At the request of Leibniz, the time was publicly extended for a year and a half. At 4 p.m. on 29 January 1697 when he arrived home from the Royal Mint, Isaac Newton found the challenge in a letter from Johann Bernoulli. Newton stayed up all night to solve it and mailed the solution anonymously by the next post. Upon reading the solution, Bernoulli immediately recognized its author, exclaiming that he "recognizes a lion from his claw mark". This story gives some idea of Newton's power, since Johann Bernoulli took two weeks to solve it.D.T. Whiteside, Newton the Mathematician, in Bechler, Contemporary Newtonian Research, p. 122. Newton also wrote, "I do not love to be dunned [pestered] and teased by foreigners about mathematical things...", and Newton had already solved Newton's minimal resistance problem, which is considered the first of the kind in calculus of variations.
In the end, five mathematicians responded with solutions: Newton, Jakob Bernoulli, Gottfried Leibniz, Ehrenfried Walther von Tschirnhaus and Guillaume de l'Hôpital. Four of the solutions (excluding l'Hôpital's) were published in the same edition of the journal as Johann Bernoulli's. In his paper, Jakob Bernoulli gave a proof of the condition for least time similar to that below before showing that its solution is a cycloid. According to Newtonian scholar Tom Whiteside, in an attempt to outdo his brother, Jakob Bernoulli created a harder version of the brachistochrone problem. In solving it, he developed new methods that were refined by Leonhard Euler into what the latter called (in 1766) the calculus of variations. Joseph-Louis Lagrange did further work that resulted in modern infinitesimal calculus.
Johann Bernoulli's solution
Introduction
In a letter to L’Hôpital, (21/12/1696), Bernoulli stated that when considering the problem of the curve of quickest descent, after only 2 days he noticed a curious affinity or connection with another no less remarkable problem leading to an ‘indirect method’ of solution. Then shortly afterwards he discovered a ‘direct method’.
Direct method
In a letter to Henri Basnage, held at the University of Basel Public Library, dated 30 March 1697, Johann Bernoulli stated that he had found two methods (always referred to as "direct" and "indirect") to show that the Brachistochrone was the "common cycloid", also called the "roulette". Following advice from Leibniz, he included only the indirect method in the Acta Eruditorum Lipsidae of May 1697. He wrote that this was partly because he believed it was sufficient to convince anyone who doubted the conclusion, partly because it also resolved two famous problems in optics that "the late Mr. Huygens" had raised in his treatise on light. In the same letter he criticised Newton for concealing his method.
In addition to his indirect method he also published the five other replies to the problem that he received.
Johann Bernoulli's direct method is historically important as a proof that the brachistochrone is the cycloid. The method is to determine the curvature of the curve at each point. All the other proofs, including Newton's (which was not revealed at the time) are based on finding the gradient at each point.
In 1718, Bernoulli explained how he solved the brachistochrone problem by his direct method.The Early Period of the Calculus of Variations, by P. Freguglia and M. Giaquinta, pp. 53–57, .
He explained that he had not published it in 1697, for reasons that no longer applied in 1718. This paper was largely ignored until 1904 when the depth of the method was first appreciated by Constantin Carathéodory, who stated that it shows that the cycloid is the only possible curve of quickest descent. According to him, the other solutions simply implied that the time of descent is stationary for the cycloid, but not necessarily the minimum possible.
Analytic solution
A body is regarded as sliding along any small circular arc Ce between the radii KC and Ke, with centre K fixed. The first stage of the proof involves finding the particular circular arc, Mm, which the body traverses in the minimum time.
The line KNC intersects AL at N, and line Kne intersects it at n, and they make a small angle CKe at K. Let NK = a, and define a variable point, C on KN extended. Of all the possible circular arcs Ce, it is required to find the arc Mm, which requires the minimum time to slide between the 2 radii, KM and Km. To find Mm Bernoulli argues as follows.
Let MN = x. He defines m so that MD = mx, and n so that Mm = nx + na and notes that x is the only variable and that m is finite and n is infinitely small. The small time to travel along arc Mm is , which has to be a minimum (‘un plus petit’). He does not explain that because Mm is so small the speed along it can be assumed to be the speed at M, which is as the square root of MD, the vertical distance of M below the horizontal line AL.
It follows that, when differentiated this must give
so that x = a.
This condition defines the curve that the body slides along in the shortest time possible. For each point, M on the curve, the radius of curvature, MK is cut in 2 equal parts by its axis AL. This property, which Bernoulli says had been known for a long time, is unique to the cycloid.
Finally, he considers the more general case where the speed is an arbitrary function X(x), so the time to be minimised is .
The minimum condition then becomes
which he writes as :
and which gives MN (=x) as a function of NK (= a). From this the equation of the curve could be obtained from the integral calculus, though he does not demonstrate this.
Synthetic solution
He then proceeds with what he called his Synthetic Solution, which was a classical, geometrical proof, that there is only a single curve that a body can slide down in the minimum time, and that curve is the cycloid.
"The reason for the synthetic demonstration, in the manner of the ancients, is to convince Mr. de la Hire. He has little time for our new analysis, describing it as false (He claims he has found 3 ways to prove that the curve is a cubic parabola)" – Letter from Johan Bernoulli to Pierre Varignon dated 27 Jul 1697.
Assume AMmB is the part of the cycloid joining A to B, which the body slides down in the minimum time. Let ICcJ be part of a different curve joining A to B, which can be closer to AL than AMmB. If the arc Mm subtends the angle MKm at its centre of curvature, K, let the arc on IJ that subtends the same angle be Cc. The circular arc through C with centre K is Ce. Point D on AL is vertically above M. Join K to D and point H is where CG intersects KD, extended if necessary.
Let and t be the times the body takes to fall along Mm and Ce respectively.
, ,
Extend CG to point F where, and since , it follows that
Since MN = NK, for the cycloid:
, , and
If Ce is closer to K than Mm then
and
In either case,
, and it follows that
If the arc, Cc subtended by the angle infinitesimal angle MKm on IJ is not circular, it must be greater than Ce, since Cec becomes a right-triangle in the limit as angle MKm approaches zero.
Note, Bernoulli proves that CF > CG by a similar but different argument.
From this he concludes that a body traverses the cycloid AMB in less time than any other curve ACB.
Indirect method
According to Fermat’s principle, the actual path between two points taken by a beam of light (which obeys Snell's law of refraction) is one that takes the least time. In 1697 Johann Bernoulli used this principle to derive the brachistochrone curve by considering the trajectory of a beam of light in a medium where the speed of light increases following a constant vertical acceleration (that of gravity g).
By the conservation of energy, the instantaneous speed of a body v after falling a height y in a uniform gravitational field is given by:
,
The speed of motion of the body along an arbitrary curve does not depend on the horizontal displacement.
Bernoulli noted that Snell's law of refraction gives a constant of the motion for a beam of light in a medium of variable density:
,
where vm is the constant and represents the angle of the trajectory with respect to the vertical.
The equations above lead to two conclusions:
At the onset, the angle must be zero when the particle speed is zero. Hence, the brachistochrone curve is tangent to the vertical at the origin.
The speed reaches a maximum value when the trajectory becomes horizontal and the angle θ = 90°.
Assuming for simplicity that the particle (or the beam) with coordinates (x,y) departs from the point (0,0) and reaches maximum speed after falling a vertical distance D:
.
Rearranging terms in the law of refraction and squaring gives:
which can be solved for dx in terms of dy:
.
Substituting from the expressions for v and vm above gives:
which is the differential equation of an inverted cycloid generated by a circle of diameter D=2r, whose parametric equation is:
where φ is a real parameter, corresponding to the angle through which the rolling circle has rotated. For given φ, the circle's centre lies at .
In the brachistochrone problem, the motion of the body is given by the time evolution of the parameter:
where t is the time since the release of the body from the point (0,0).
Jakob Bernoulli's solution
Johann's brother Jakob showed how 2nd differentials can be used to obtain the condition for least time. A modernized version of the proof is as follows. If we make a negligible deviation from the path of least time, then, for the differential triangle formed by the displacement along the path and the horizontal and vertical displacements,
.
On differentiation with dy fixed we get,
.
And finally rearranging terms gives,
where the last part is the displacement for given change in time for 2nd differentials. Now consider the changes along the two neighboring paths in the figure below for which the horizontal separation between paths along the central line is d2x (the same for both the upper and lower differential triangles). Along the old and new paths, the parts that differ are,
For the path of least times these times are equal so for their difference we get,
And the condition for least time is,
which agrees with Johann's assumption based on the law of refraction.
Newton's solution
Introduction
In June 1696, Johann Bernoulli had used the pages of the Acta Eruditorum Lipsidae to pose a challenge to the international mathematical community: to find the form of the curve joining two fixed points so that a mass will slide down along it, under the influence of gravity alone, in the minimum amount of time. The solution was originally to be submitted within six months. At the suggestion of Leibniz, Bernoulli extended the challenge until Easter 1697, by means of a printed text called "Programma", published in Groningen, in the Netherlands.
The Programma is dated 1 January 1697, in the Gregorian Calendar. This was 22 December 1696 in the Julian Calendar, in use in Britain.
According to Newton's niece, Catherine Conduitt, Newton learned of the challenge at 4 pm on 29 January and had solved it by 4 am the following morning. His solution, communicated to the Royal Society, is dated 30 January. This solution, later published anonymously in the Philosophical Transactions, is correct but does not indicate the method by which Newton arrived at his conclusion. Bernoulli, writing to Henri Basnage in March 1697, indicated that even though its author, "by an excess of modesty", had not revealed his name, yet even from the scant details supplied it could be recognised as Newton's work, "as the lion by its claw" (in Latin, ex ungue Leonem).
D. T. Whiteside notes that the letter in French has ex ungue Leonem preceded by the French word comme. The much quoted version tanquam ex ungue Leonem is due to David Brewster's 1855 book on the life and works of Newton. Bernoulli's intention was, Whiteside argues, simply to indicate he could tell the anonymous solution was Newton's, just as it was possible to tell that an animal was a lion given its claw; it was not meant to suggest that Bernoulli considered Newton to be the lion among mathematicians, as it has since come to be interpreted.
John Wallis, who was 80 years old at the time, had learned of the problem in September 1696 from Johann Bernoulli's youngest brother Hieronymus, and had spent three months attempting a solution before passing it in December to David Gregory, who also failed to solve it. After Newton had submitted his solution, Gregory asked him for the details and made notes from their conversation. These can be found in the University of Edinburgh Library, manuscript A , dated 7 March 1697. Either Gregory did not understand Newton's argument, or Newton's explanation was very brief. However, it is possible, with a high degree of confidence, to construct Newton's proof from Gregory's notes, by analogy with his method to determine the solid of minimum resistance (Principia, Book 2, Proposition 34, Scholium 2). A detailed description of his solution of this latter problem is included in the draft of a letter in 1694, also to David Gregory. In addition to the minimum time curve problem, there was a second problem that Newton also solved at the same time. Both solutions appeared anonymously in Philosophical Transactions of the Royal Society, for January 1697.
The Brachistochrone problem
Fig. 1, shows Gregory’s diagram (except the additional line IF is absent from it, and Z, the start point has been added). The curve ZVA is a cycloid and CHV is its generating circle. Since it appears that the body is moving upward from e to E, it must be assumed that a small body is released from Z and slides along the curve to A, without friction, under the action of gravity.
Consider a small arc eE, which the body is ascending. Assume that it traverses the straight line eL to point L, horizontally displaced from E by a small distance, o, instead of the arc eE. Note, that eL is not the tangent at e, and that o is negative when L is between B and E. Draw the line through E parallel to CH, cutting eL at n. From a property of the cycloid, En is the normal to the tangent at E, and similarly the tangent at E is parallel to VH.
Since the displacement EL is small, it differs little in direction from the tangent at E so that the angle EnL is close to a right-angle. In the limit as the arc eE approaches zero, eL becomes parallel to VH, provided o is small compared to eE making the triangles EnL and CHV similar.
Also en approaches the length of chord eE, and the increase in length, , ignoring terms in and higher, which represent the error due to the approximation that eL and VH are parallel.
The speed along eE or eL can be taken as that at E, proportional to , which is as CH, since
This appears to be all that Gregory’s note contains.
Let t be the additional time to reach L,
Therefore, the increase in time to traverse a small arc displaced at one endpoint depends only on the displacement at the endpoint and is independent of the position of the arc. However, by Newton’s method, this is just the condition required for the curve to be traversed in the minimum time possible. Therefore, he concludes that the minimum curve must be the cycloid.
He argues as follows.
Assuming now that Fig. 1 is the minimum curve not yet determined, with vertical axis CV, and the circle CHV removed, and Fig. 2 shows part of the curve between the infinitesimal arc eE and a further infinitesimal arc Ff a finite distance along the curve. The extra time, t, to traverse eL (rather than eE) is nL divided by the speed at E (proportional to ), ignoring terms in and higher:
,
At L the particle continues along a path LM, parallel to the original EF, to some arbitrary point M. As it has the same speed at L as at E, the time to traverse LM is the same as it would have been along the original curve EF. At M it returns to the original path at point f. By the same reasoning, the reduction in time, T, to reach f from M rather than from F is
The difference (t – T) is the extra time it takes along the path compared to the original :
plus terms in and higher (1)
Because is the minimum curve, (t – T) is must be greater than zero, whether o is positive or negative. It follows that the coefficient of o in (1) must be zero:
(2) in the limit as eE and fF approach zero. Note since is the minimum curve it has to be assumed that the coefficient of is greater than zero.
Clearly there has to be 2 equal and opposite displacements, or the body would not return to the endpoint, A, of the curve.
If e is fixed, and if f is considered a variable point higher up the curve, then for all such points, f, is constant (equal to ). By keeping f fixed and making e variable it is clear that is also constant.
But, since points, e and f are arbitrary, equation (2) can be true only if , everywhere, and this condition characterises the curve that is sought. This is the same technique he uses to find the form of the Solid of Least Resistance.
For the cycloid, , so that , which was shown above to be constant, and the Brachistochrone is the cycloid.
Newton gives no indication of how he discovered that the cycloid satisfied this last relation. It may have been by trial and error, or he may have recognised immediately that it implied the curve was the cycloid.
| Mathematics | Two-dimensional space | null |
171882 | https://en.wikipedia.org/wiki/Calculus%20of%20variations | Calculus of variations | The calculus of variations (or variational calculus) is a field of mathematical analysis that uses variations, which are small changes in functions
and functionals, to find maxima and minima of functionals: mappings from a set of functions to the real numbers. Functionals are often expressed as definite integrals involving functions and their derivatives. Functions that maximize or minimize functionals may be found using the Euler–Lagrange equation of the calculus of variations.
A simple example of such a problem is to find the curve of shortest length connecting two points. If there are no constraints, the solution is a straight line between the points. However, if the curve is constrained to lie on a surface in space, then the solution is less obvious, and possibly many solutions may exist. Such solutions are known as geodesics. A related problem is posed by Fermat's principle: light follows the path of shortest optical length connecting two points, which depends upon the material of the medium. One corresponding concept in mechanics is the principle of least/stationary action.
Many important problems involve functions of several variables. Solutions of boundary value problems for the Laplace equation satisfy the Dirichlet's principle. Plateau's problem requires finding a surface of minimal area that spans a given contour in space: a solution can often be found by dipping a frame in soapy water. Although such experiments are relatively easy to perform, their mathematical formulation is far from simple: there may be more than one locally minimizing surface, and they may have non-trivial topology.
History
The calculus of variations may be said to begin with Newton's minimal resistance problem in 1687, followed by the brachistochrone curve problem raised by Johann Bernoulli (1696). It immediately occupied the attention of Jacob Bernoulli and the Marquis de l'Hôpital, but Leonhard Euler first elaborated the subject, beginning in 1733. Joseph-Louis Lagrange was influenced by Euler's work to contribute significantly to the theory. After Euler saw the 1755 work of the 19-year-old Lagrange, Euler dropped his own partly geometric approach in favor of Lagrange's purely analytic approach and renamed the subject the calculus of variations in his 1756 lecture Elementa Calculi Variationum.
Adrien-Marie Legendre (1786) laid down a method, not entirely satisfactory, for the discrimination of maxima and minima. Isaac Newton and Gottfried Leibniz also gave some early attention to the subject. To this discrimination Vincenzo Brunacci (1810), Carl Friedrich Gauss (1829), Siméon Poisson (1831), Mikhail Ostrogradsky (1834), and Carl Jacobi (1837) have been among the contributors. An important general work is that of Pierre Frédéric Sarrus (1842) which was condensed and improved by Augustin-Louis Cauchy (1844). Other valuable treatises and memoirs have been written by Strauch (1849), John Hewitt Jellett (1850), Otto Hesse (1857), Alfred Clebsch (1858), and Lewis Buffett Carll (1885), but perhaps the most important work of the century is that of Karl Weierstrass. His celebrated course on the theory is epoch-making, and it may be asserted that he was the first to place it on a firm and unquestionable foundation. The 20th and the 23rd Hilbert problem published in 1900 encouraged further development.
In the 20th century David Hilbert, Oskar Bolza, Gilbert Ames Bliss, Emmy Noether, Leonida Tonelli, Henri Lebesgue and Jacques Hadamard among others made significant contributions. Marston Morse applied calculus of variations in what is now called Morse theory. Lev Pontryagin, Ralph Rockafellar and F. H. Clarke developed new mathematical tools for the calculus of variations in optimal control theory. The dynamic programming of Richard Bellman is an alternative to the calculus of variations.
Extrema
The calculus of variations is concerned with the maxima or minima (collectively called extrema) of functionals. A functional maps functions to scalars, so functionals have been described as "functions of functions." Functionals have extrema with respect to the elements of a given function space defined over a given domain. A functional is said to have an extremum at the function if has the same sign for all in an arbitrarily small neighborhood of The function is called an extremal function or extremal. The extremum is called a local maximum if everywhere in an arbitrarily small neighborhood of and a local minimum if there. For a function space of continuous functions, extrema of corresponding functionals are called strong extrema or weak extrema, depending on whether the first derivatives of the continuous functions are respectively all continuous or not.
Both strong and weak extrema of functionals are for a space of continuous functions but strong extrema have the additional requirement that the first derivatives of the functions in the space be continuous. Thus a strong extremum is also a weak extremum, but the converse may not hold. Finding strong extrema is more difficult than finding weak extrema. An example of a necessary condition that is used for finding weak extrema is the Euler–Lagrange equation.
Euler–Lagrange equation
Finding the extrema of functionals is similar to finding the maxima and minima of functions. The maxima and minima of a function may be located by finding the points where its derivative vanishes (i.e., is equal to zero). The extrema of functionals may be obtained by finding functions for which the functional derivative is equal to zero. This leads to solving the associated Euler–Lagrange equation.
Consider the functional
where
are constants,
is twice continuously differentiable,
is twice continuously differentiable with respect to its arguments and
If the functional attains a local minimum at and is an arbitrary function that has at least one derivative and vanishes at the endpoints and then for any number close to 0,
The term is called the variation of the function and is denoted by
Substituting for in the functional the result is a function of
Since the functional has a minimum for the function has a minimum at and thus,
Taking the total derivative of where and are considered as functions of rather than yields
and because and
Therefore,
where when and we have used integration by parts on the second term. The second term on the second line vanishes because at and by definition. Also, as previously mentioned the left side of the equation is zero so that
According to the fundamental lemma of calculus of variations, the part of the integrand in parentheses is zero, i.e.
which is called the Euler–Lagrange equation. The left hand side of this equation is called the functional derivative of and is denoted or
In general this gives a second-order ordinary differential equation which can be solved to obtain the extremal function The Euler–Lagrange equation is a necessary, but not sufficient, condition for an extremum A sufficient condition for a minimum is given in the section Variations and sufficient condition for a minimum.
Example
In order to illustrate this process, consider the problem of finding the extremal function which is the shortest curve that connects two points and The arc length of the curve is given by
with
Note that assuming is a function of loses generality; ideally both should be a function of some other parameter. This approach is good solely for instructive purposes.
The Euler–Lagrange equation will now be used to find the extremal function that minimizes the functional
with
Since does not appear explicitly in the first term in the Euler–Lagrange equation vanishes for all and thus,
Substituting for and taking the derivative,
Thus
for some constant Then
where
Solving, we get
which implies that
is a constant and therefore that the shortest curve that connects two points and is
and we have thus found the extremal function that minimizes the functional so that is a minimum. The equation for a straight line is In other words, the shortest distance between two points is a straight line.
Beltrami's identity
In physics problems it may be the case that meaning the integrand is a function of and but does not appear separately. In that case, the Euler–Lagrange equation can be simplified to the Beltrami identity
where is a constant. The left hand side is the Legendre transformation of with respect to
The intuition behind this result is that, if the variable is actually time, then the statement implies that the Lagrangian is time-independent. By Noether's theorem, there is an associated conserved quantity. In this case, this quantity is the Hamiltonian, the Legendre transform of the Lagrangian, which (often) coincides with the energy of the system. This is (minus) the constant in Beltrami's identity.
Euler–Poisson equation
If depends on higher-derivatives of that is, if then must satisfy the Euler–Poisson equation,
Du Bois-Reymond's theorem
The discussion thus far has assumed that extremal functions possess two continuous derivatives, although the existence of the integral requires only first derivatives of trial functions. The condition that the first variation vanishes at an extremal may be regarded as a weak form of the Euler–Lagrange equation. The theorem of Du Bois-Reymond asserts that this weak form implies the strong form. If has continuous first and second derivatives with respect to all of its arguments, and if
then has two continuous derivatives, and it satisfies the Euler–Lagrange equation.
Lavrentiev phenomenon
Hilbert was the first to give good conditions for the Euler–Lagrange equations to give a stationary solution. Within a convex area and a positive thrice differentiable Lagrangian the solutions are composed of a countable collection of sections that either go along the boundary or satisfy the Euler–Lagrange equations in the interior.
However Lavrentiev in 1926 showed that there are circumstances where there is no optimum solution but one can be approached arbitrarily closely by increasing numbers of sections. The Lavrentiev Phenomenon identifies a difference in the infimum of a minimization problem across different classes of admissible functions. For instance the following problem, presented by Manià in 1934:
Clearly, minimizes the functional, but we find any function gives a value bounded away from the infimum.
Examples (in one-dimension) are traditionally manifested across and but Ball and Mizel procured the first functional that displayed Lavrentiev's Phenomenon across and for There are several results that gives criteria under which the phenomenon does not occur - for instance 'standard growth', a Lagrangian with no dependence on the second variable, or an approximating sequence satisfying Cesari's Condition (D) - but results are often particular, and applicable to a small class of functionals.
Connected with the Lavrentiev Phenomenon is the repulsion property: any functional displaying Lavrentiev's Phenomenon will display the weak repulsion property.
Functions of several variables
For example, if denotes the displacement of a membrane above the domain in the plane, then its potential energy is proportional to its surface area:
Plateau's problem consists of finding a function that minimizes the surface area while assuming prescribed values on the boundary of ; the solutions are called minimal surfaces. The Euler–Lagrange equation for this problem is nonlinear:
See Courant (1950) for details.
Dirichlet's principle
It is often sufficient to consider only small displacements of the membrane, whose energy difference from no displacement is approximated by
The functional is to be minimized among all trial functions that assume prescribed values on the boundary of If is the minimizing function and is an arbitrary smooth function that vanishes on the boundary of then the first variation of must vanish:
Provided that u has two derivatives, we may apply the divergence theorem to obtain
where is the boundary of is arclength along and is the normal derivative of on Since vanishes on and the first variation vanishes, the result is
for all smooth functions that vanish on the boundary of The proof for the case of one dimensional integrals may be adapted to this case to show that
in
The difficulty with this reasoning is the assumption that the minimizing function must have two derivatives. Riemann argued that the existence of a smooth minimizing function was assured by the connection with the physical problem: membranes do indeed assume configurations with minimal potential energy. Riemann named this idea the Dirichlet principle in honor of his teacher Peter Gustav Lejeune Dirichlet. However Weierstrass gave an example of a variational problem with no solution: minimize
among all functions that satisfy and
can be made arbitrarily small by choosing piecewise linear functions that make a transition between −1 and 1 in a small neighborhood of the origin. However, there is no function that makes Eventually it was shown that Dirichlet's principle is valid, but it requires a sophisticated application of the regularity theory for elliptic partial differential equations; see Jost and Li–Jost (1998).
Generalization to other boundary value problems
A more general expression for the potential energy of a membrane is
This corresponds to an external force density in an external force on the boundary and elastic forces with modulus acting on The function that minimizes the potential energy with no restriction on its boundary values will be denoted by Provided that and are continuous, regularity theory implies that the minimizing function will have two derivatives. In taking the first variation, no boundary condition need be imposed on the increment The first variation of is given by
If we apply the divergence theorem, the result is
If we first set on the boundary integral vanishes, and we conclude as before that
in Then if we allow to assume arbitrary boundary values, this implies that must satisfy the boundary condition
on This boundary condition is a consequence of the minimizing property of : it is not imposed beforehand. Such conditions are called natural boundary conditions.
The preceding reasoning is not valid if vanishes identically on In such a case, we could allow a trial function where is a constant. For such a trial function,
By appropriate choice of can assume any value unless the quantity inside the brackets vanishes. Therefore, the variational problem is meaningless unless
This condition implies that net external forces on the system are in equilibrium. If these forces are in equilibrium, then the variational problem has a solution, but it is not unique, since an arbitrary constant may be added. Further details and examples are in Courant and Hilbert (1953).
Eigenvalue problems
Both one-dimensional and multi-dimensional eigenvalue problems can be formulated as variational problems.
Sturm–Liouville problems
The Sturm–Liouville eigenvalue problem involves a general quadratic form
where is restricted to functions that satisfy the boundary conditions
Let be a normalization integral
The functions and are required to be everywhere positive and bounded away from zero. The primary variational problem is to minimize the ratio among all satisfying the endpoint conditions, which is equivalent to minimizing under the constraint that is constant. It is shown below that the Euler–Lagrange equation for the minimizing is
where is the quotient
It can be shown (see Gelfand and Fomin 1963) that the minimizing has two derivatives and satisfies the Euler–Lagrange equation. The associated will be denoted by ; it is the lowest eigenvalue for this equation and boundary conditions. The associated minimizing function will be denoted by This variational characterization of eigenvalues leads to the Rayleigh–Ritz method: choose an approximating as a linear combination of basis functions (for example trigonometric functions) and carry out a finite-dimensional minimization among such linear combinations. This method is often surprisingly accurate.
The next smallest eigenvalue and eigenfunction can be obtained by minimizing under the additional constraint
This procedure can be extended to obtain the complete sequence of eigenvalues and eigenfunctions for the problem.
The variational problem also applies to more general boundary conditions. Instead of requiring that vanish at the endpoints, we may not impose any condition at the endpoints, and set
where and are arbitrary. If we set , the first variation for the ratio is
where λ is given by the ratio as previously.
After integration by parts,
If we first require that vanish at the endpoints, the first variation will vanish for all such only if
If satisfies this condition, then the first variation will vanish for arbitrary only if
These latter conditions are the natural boundary conditions for this problem, since they are not imposed on trial functions for the minimization, but are instead a consequence of the minimization.
Eigenvalue problems in several dimensions
Eigenvalue problems in higher dimensions are defined in analogy with the one-dimensional case. For example, given a domain with boundary in three dimensions we may define
and
Let be the function that minimizes the quotient
with no condition prescribed on the boundary The Euler–Lagrange equation satisfied by is
where
The minimizing must also satisfy the natural boundary condition
on the boundary This result depends upon the regularity theory for elliptic partial differential equations; see Jost and Li–Jost (1998) for details. Many extensions, including completeness results, asymptotic properties of the eigenvalues and results concerning the nodes of the eigenfunctions are in Courant and Hilbert (1953).
Applications
Optics
Fermat's principle states that light takes a path that (locally) minimizes the optical length between its endpoints. If the -coordinate is chosen as the parameter along the path, and along the path, then the optical length is given by
where the refractive index depends upon the material.
If we try then the first variation of (the derivative of with respect to ε) is
After integration by parts of the first term within brackets, we obtain the Euler–Lagrange equation
The light rays may be determined by integrating this equation. This formalism is used in the context of Lagrangian optics and Hamiltonian optics.
Snell's law
There is a discontinuity of the refractive index when light enters or leaves a lens. Let
where and are constants. Then the Euler–Lagrange equation holds as before in the region where or and in fact the path is a straight line there, since the refractive index is constant. At the must be continuous, but may be discontinuous. After integration by parts in the separate regions and using the Euler–Lagrange equations, the first variation takes the form
The factor multiplying is the sine of angle of the incident ray with the axis, and the factor multiplying is the sine of angle of the refracted ray with the axis. Snell's law for refraction requires that these terms be equal. As this calculation demonstrates, Snell's law is equivalent to vanishing of the first variation of the optical path length.
Fermat's principle in three dimensions
It is expedient to use vector notation: let let be a parameter, let be the parametric representation of a curve and let be its tangent vector. The optical length of the curve is given by
Note that this integral is invariant with respect to changes in the parametric representation of The Euler–Lagrange equations for a minimizing curve have the symmetric form
where
It follows from the definition that satisfies
Therefore, the integral may also be written as
This form suggests that if we can find a function whose gradient is given by then the integral is given by the difference of at the endpoints of the interval of integration. Thus the problem of studying the curves that make the integral stationary can be related to the study of the level surfaces of In order to find such a function, we turn to the wave equation, which governs the propagation of light. This formalism is used in the context of Lagrangian optics and Hamiltonian optics.
Connection with the wave equation
The wave equation for an inhomogeneous medium is
where is the velocity, which generally depends upon Wave fronts for light are characteristic surfaces for this partial differential equation: they satisfy
We may look for solutions in the form
In that case, satisfies
where According to the theory of first-order partial differential equations, if then satisfies
along a system of curves (the light rays) that are given by
These equations for solution of a first-order partial differential equation are identical to the Euler–Lagrange equations if we make the identification
We conclude that the function is the value of the minimizing integral as a function of the upper end point. That is, when a family of minimizing curves is constructed, the values of the optical length satisfy the characteristic equation corresponding the wave equation. Hence, solving the associated partial differential equation of first order is equivalent to finding families of solutions of the variational problem. This is the essential content of the Hamilton–Jacobi theory, which applies to more general variational problems.
Mechanics
In classical mechanics, the action, is defined as the time integral of the Lagrangian, The Lagrangian is the difference of energies,
where is the kinetic energy of a mechanical system and its potential energy. Hamilton's principle (or the action principle) states that the motion of a conservative holonomic (integrable constraints) mechanical system is such that the action integral
is stationary with respect to variations in the path
The Euler–Lagrange equations for this system are known as Lagrange's equations:
and they are equivalent to Newton's equations of motion (for such systems).
The conjugate momenta are defined by
For example, if
then
Hamiltonian mechanics results if the conjugate momenta are introduced in place of by a Legendre transformation of the Lagrangian into the Hamiltonian defined by
The Hamiltonian is the total energy of the system:
Analogy with Fermat's principle suggests that solutions of Lagrange's equations (the particle trajectories) may be described in terms of level surfaces of some function of This function is a solution of the Hamilton–Jacobi equation:
Further applications
Further applications of the calculus of variations include the following:
The derivation of the catenary shape
Solution to Newton's minimal resistance problem
Solution to the brachistochrone problem
Solution to the tautochrone problem
Solution to isoperimetric problems
Calculating geodesics
Finding minimal surfaces and solving Plateau's problem
Optimal control
Analytical mechanics, or reformulations of Newton's laws of motion, most notably Lagrangian and Hamiltonian mechanics;
Geometric optics, especially Lagrangian and Hamiltonian optics;
Variational method (quantum mechanics), one way of finding approximations to the lowest energy eigenstate or ground state, and some excited states;
Variational Bayesian methods, a family of techniques for approximating intractable integrals arising in Bayesian inference and machine learning;
Variational methods in general relativity, a family of techniques using calculus of variations to solve problems in Einstein's general theory of relativity;
Finite element method is a variational method for finding numerical solutions to boundary-value problems in differential equations;
Total variation denoising, an image processing method for filtering high variance or noisy signals.
Variations and sufficient condition for a minimum
Calculus of variations is concerned with variations of functionals, which are small changes in the functional's value due to small changes in the function that is its argument. The first variation is defined as the linear part of the change in the functional, and the second variation is defined as the quadratic part.
For example, if is a functional with the function as its argument, and there is a small change in its argument from to where is a function in the same function space as then the corresponding change in the functional is
The functional is said to be differentiable if
where is a linear functional, is the norm of and as The linear functional is the first variation of and is denoted by,
The functional is said to be twice differentiable if
where is a linear functional (the first variation), is a quadratic functional, and as The quadratic functional is the second variation of and is denoted by,
The second variation is said to be strongly positive if
for all and for some constant .
Using the above definitions, especially the definitions of first variation, second variation, and strongly positive, the following sufficient condition for a minimum of a functional can be stated.
| Mathematics | Calculus and analysis | null |
171885 | https://en.wikipedia.org/wiki/Family%20farm | Family farm | A family farm is generally understood to be a farm owned and/or operated by a family. It is sometimes considered to be an estate passed down by inheritance.
Although a recurring conceptual and archetypal distinction is that of a family farm as a smallholding versus corporate farming as large-scale agribusiness, that notion does not accurately describe the realities of farm ownership in many countries. Family farm businesses can take many forms, from smallholder farms to larger farms operated under intensive farming practices. In various countries, most farm families have structured their farm businesses as corporations (such as limited liability companies) or trusts, for liability, tax, and business purposes. Thus, the idea of a family farm as a unitary concept or definition does not easily translate across languages, cultures, or centuries, as there are substantial differences in agricultural traditions and histories between countries and between centuries within a country. For example, in U.S. agriculture, a family farm can be of any size, as long as the ownership is held within a family. A 2014 USDA report shows that family farms operate 90 percent of the nation's farmland, and account for 85 percent of the country's agricultural production value. However, that does not at all imply that corporate farming is a small presence in U.S. agriculture; rather, it simply reflects the fact that many corporations are closely held. In contrast, in Brazilian agriculture, the official definition of a family farm (agricultura familiar) is limited to small farms worked primarily by members of a single family; but again, this fact does not imply that corporate farming is a small presence in Brazilian agriculture; rather, it simply reflects the fact that large farms with many workers cannot be legally classified under the family farm label because that label is legally reserved for smallholdings in that country.
Farms that would not be considered family farms would be those operated as collectives, non-family corporations, or in other institutionalised forms. At least 500 million of the world's [estimated] 570 million farms are managed by families, making family farms predominant in global agriculture.
Definitions
An "informal discussion of the concepts and definitions" in a working paper published by Food and Agriculture Organization of the United Nations in 2014 reviewed English, Spanish and French definitions of the concept of "family farm".
Definitions referred to one or more of labor, management, size, provision of family livelihood, residence, family ties and generational aspects, community and social networks, subsistence orientation, patrimony, land ownership and family investment. The disparity of definitions reflects national and geographical differences in cultures, rural land tenure, and rural economies, as well as the different purposes for which definitions are coined.
The 2012 United States Census of Agriculture defines a family farm as "any farm where the majority of the business is owned by the operator and individuals related to the operator, including relatives who do not live in the operator’s household"; it defines a farm as "any place from which $1,000 or more of agricultural products were produced and sold, or normally would have been sold, during a given year."
The Food and Agriculture Organization of the United Nations defines a "family farm" as one that relies primarily on family members for labour and management.
In some usages, "family farm" implies that the farm remains within the ownership of a family over a number of generations.
Being special-purpose definitions, the definitions found in laws or regulations may differ substantially from commonly understood meanings of "family farm". For example, In the United States, under federal Farm Ownership loan regulations, the definition of a "family farm" does not specify the nature of farm ownership, and management of the farm is either by the borrower, or by members operating the farm when a loan is made to a corporation, co-operative or other entity. The complete definition can be found in the US Code of Federal Regulations 7 CFR 1943.4.
History
In the Roman Republic, latifundia, great landed estates, specialised in agriculture destined for export, producing grain, olive oil, or wine, corresponding largely to modern industrialized agriculture but depending on slave labour instead of mechanization, developed after the Second Punic War and increasingly replaced the former system of family-owned small or intermediate farms in the Roman Empire period. The basis of the latifundia in Spain and Sicily was the ager publicus that fell to the dispensation of the state through Rome's policy of war in the 1st century BC and the 1st century AD.
In the collapse of the Western Roman Empire, the largely self-sufficient villa-system of the latifundia remained among the few political-cultural centres of a fragmented Europe. These latifundia had been of great importance economically, until the long-distance shipping of wine and oil, grain and garum disintegrated, but extensive lands controlled in a single pair of hands still constituted power: it can be argued that the latifundia formed part of the economic basis of the European social feudal system, taking the form of Manorialism, the essential element of feudal society, and the organizing principle of rural economy in medieval Europe.
Manorialism was characterised by the vesting of legal and economic power in a Lord of the Manor, supported economically from his own direct landholding in a manor (sometimes called a fief), and from the obligatory contributions of a legally subject part of the peasant population under the jurisdiction of himself and his manorial court.
Manorialism died slowly and piecemeal, along with its most vivid feature in the landscape, the open field system.
It outlasted serfdom as it outlasted feudalism: "primarily an economic organization, it could maintain a warrior, but it could equally well maintain a capitalist landlord. It could be self-sufficient, yield produce for the market, or it could yield a money rent." The last feudal dues in France were abolished at the French Revolution. In parts of eastern Germany, the Rittergut manors of Junkers remained until World War II.
The common law of the leasehold estate relation evolved in medieval England. That law still retains many archaic terms and principles pertinent to a feudal social order. Under the tenant system, a farm may be worked by the same family over many generations, but what is inherited is not the farm's estate itself but the lease on the estate.
In much of Europe, serfdom was abolished only in the modern period, in Western Europe after the French Revolution, in Russia as late as in 1861.
In contrast to the Roman system of latifundia and the derived system of manoralism, the Germanic peoples had a system based on heritable estates owned by individual families or clans.
The Germanic term for "heritable estate, allodium" was *ōþalan (Old English ēþel), which incidentally was also used as a rune name; the gnomic verse on this term in the Anglo-Saxon rune poem reads:
"[An estate] is very dear to every man, if he can enjoy there in his house whatever is right and proper in constant prosperity."
In the inheritance system known as Salic patrimony (also gavelkind in its exceptional survival in medieval Kent)
refers to this clan-based possession of real estate property, particularly in Germanic context. Terra salica could not be sold or otherwise disposed; it was not alienable.
Much of Germanic Europe has a history of overlap or conflict between the feudal system of manoralism, where the estate is owned by noblemen and leased to the tenants or worked by serfs, and the Germanic system of free farmers working landed estates heritable within their clan or family. Historical prevalence of the Germanic system of independent estates or Höfe resulted in dispersed settlement (Streusiedlung) structure, as opposed to the village-centered settlements of manoralism.
In German-speaking Europe, a farmyard is known as a Hof; in modern German this word designates the area enclosed by the farm buildings, not the fields around them, and it is also used in other everyday situations for courtyards of any type (Hinterhof = 'back yard', etc.). The recharacterized compound Bauernhof was formed in the early modern period to designate family farming estates and today is the most common word for 'farm', while the archaic Meierhof designated a manorial estate. Historically, the unmarked term Hof was increasingly used for the royal or noble court.
The estate as a whole is referred to by the collective Gehöft (15th century); the corresponding Slavic concept being Khutor.
Höfeordnung is the German legal term for the inheritance laws regarding family farms, deriving from inheritance under medieval Saxon law.
In England, the title of yeoman was applied to such land-owning commoners from the 15th century.
In the early modern and modern period, the dissolution of manoralism went parallel to the development of intensive farming parallel to the Industrial Revolution. Mechanization enabled the cultivation of much larger areas than what was typical for the traditional estates aimed at subsistence farming, resulting in the emergence of a smaller number of large farms, with the displaced population partly contributing to the new class of industrial wage-labourers and partly emigrating to the New World or the Russian Empire (following the 1861 emancipation of the serfs). The family farms established in Imperial Russia were again collectivized under the Soviet Union, but the emigration of European farmers displaced by the Industrial Revolution contributed to the emergence of a system of family estates in the Americas (Homestead Act of 1862).
Thomas Jefferson's argument that a large number of family estates are a factor in ensuring the stability of democracy was repeatedly used in support of subsidies.
Developed world
Perceptions of the family farm
In developed countries the family farm is viewed sentimentally, as a lifestyle to be preserved for tradition's sake, or as a birthright. It is in these nations very often a political rallying cry against change in agricultural policy, most commonly in France, Japan, and the United States, where rural lifestyles are often regarded as desirable. In these countries, strange bedfellows can often be found arguing for similar measures despite otherwise vast differences in political ideology. For example, Pat Buchanan and Ralph Nader, both candidates for the office of President of the United States, held rural rallies together and spoke for measures to preserve the so-called family farm. On other economic matters they were seen as generally opposed, but found common ground on this one.
The social roles of family farms are much changed today. Until recently, staying in line with traditional and conservative sociology, the heads of the household were usually the oldest man followed closely by his oldest sons. The wife generally took care of the housework, child rearing, and financial matters pertaining to the farm. However, agricultural activities have taken on many forms and change over time. Agronomy, horticulture, aquaculture, silviculture, and apiculture, along with traditional plants and animals, all make up aspects of today's family farm. Farm wives often need to find work away from the farm to supplement farm income and children sometimes have no interest in farming as their chosen field of work.
Bolder promoters argue that as agriculture has become more efficient with the application of modern management and new technologies in each generation, the idealized classic family farm is now simply obsolete, or more often, unable to compete without the economies of scale available to larger and more modern farms. Advocates argue that family farms in all nations need to be protected, as the basis of rural society and social stability.
Viability
According to the United States Department of Agriculture, ninety-eight percent of all farms in the U.S. are family farms. Two percent of farms are not family farms, and those two percent make up fourteen percent of total agricultural output in the United States, although half of them have total sales of less than $50,000 per year. Overall, ninety-one percent of farms in the United States are considered "small family farms" (with sales of less than $250,000 per year), and those farms produce twenty-seven percent of U.S. agricultural output.
Depending on the type and size of independently owned operation, some limiting factors are:
Economies of scale: Larger farms are able to bargain more competitively, purchase more competitively, profit from economic highs, and weather lows more readily through monetary inertia than smaller farms.
Cost of inputs: fertilizer and other agrichemicals can fluctuate dramatically from season to season, partially based on oil prices, a range of 25% to 200% is common over a period of a few years.
oil prices: Directly (for farm machinery) and somewhat less directly (long-distance transport; production cost of agrichemicals), the cost of oil significantly impacts the year-to-year viability of all mechanized conventional farms.
commodity futures: the predicted price of commodity crops, hogs, grain, etc., can determine ahead of a season what seems economically viable to grow.
technology user agreements: a less publicly known factor, patented GE seed that is widely used for many crops, like cotton and soy, comes with restrictions on use, which can even include who the crop can be sold to.
wholesale infrastructure: A farmer growing larger quantities of a crop than can be sold directly to consumers has to meet a range of criteria for sale into the wholesale market, which include harvest timing and graded quality, and may also include variety, therefore, the market channel really determines most aspects of the farm decisionmaking.
availability of financing: Larger farms today often rely on lines of credit, typically from banks, to purchase the agrichemicals, and other supplies needed for each growing year. These lines are heavily affected by almost all of the other constraining factors.
government economic intervention: In some countries, notably the US and EU, government subsidies to farmers, intended to mitigate the impact on domestic farmers of economic and political activities in other areas of the economy, can be a significant source of farm income. Bailouts, when crises such as drought or the "mad cow disease" problems hit agricultural sectors, are also relied on. To some large degree, this situation is a result of the large-scale global markets farms have no alternative but to participate in.
government and industry regulation: A wide range of quotas, marketing boards and legislation governing agriculture impose complicated limits, and often require significant resources to navigate. For example, on the small farming end, in many jurisdictions, there are severe limits or prohibitions on the sale of livestock, dairy and eggs. These have arisen from pressures from all sides: food safety, environmental, industry marketing.
real estate prices: The growth of urban centers around the world, and the resulting urban sprawl have caused the price of centrally located farmland to skyrocket, while reducing the local infrastructure necessary to support farming, putting effectively intense pressure on many farmers to sell out.
Over the 20th century, the people of developed nations have collectively taken most of the steps down the path to this situation. Individual farmers opted for successive waves of new technology, happily "trading in their horses for a tractor", increasing their debt and their production capacity. This in turn required larger, more distant markets, and heavier and more complex financing. The public willingly purchased increasingly commoditized, processed, shipped and relatively inexpensive food. The availability of an increasingly diverse supply of fresh, uncured, unpreserved produce and meat in all seasons of the year (oranges in January, freshly killed steers in July, fresh pork rather than salted, smoked, or potassium-impregnated ham) opened an entirely new cuisine and an unprecedented healthy diet to millions of consumers who had never enjoyed such produce before. These abilities also brought to market an unprecedented variety of processed foods, such as corn syrup and bleached flour. For the family farm this new technology and increasingly complex marketing strategy has presented new and unprecedented challenges, and not all family farmers have been able to effectively cope with the changing market conditions.
Local food and the organic movement
In the last few decades there has been a resurgence of interest in organic and free range foods. A percentage of consumers have begun to question the viability of industrial agriculture practices and have turned to organic groceries that sell products produced on family farms including not only meat and produce but also such things as wheat germ breads and natural lye soaps (as opposed to bleached white breads and petroleum based detergent bars). Others buy these products direct from family farms. The "new family farm" provides an alternative market in some localities with an array of traditionally and naturally produced products.
Such "organic" and "free-range" farming is attainable where a significant number of affluent urban and suburban consumers willingly pay a premium for the ideals of "locally produced produce" and "humane treatment of animals". Sometimes, these farms are hobby or part-time ventures, or supported by wealth from other sources. Viable farms on a scale sufficient to support modern families at an income level commensurate with urban and suburban upper-middle-class families are often large scale operations, both in area and capital requirements. These farms, family owned and operated in a technologically and economically conventional manner, produce crops and animal products oriented to national and international markets, rather than to local markets. In assessing this complex economic situation, it is important to consider all sources of income available to these farms; for instance, the millions of dollars in farm subsidies which the United States government offers each year. As fuel prices rise, foods shipped to national and international markets are already rising in price.
United States
In 2012, the United States had 2,039,093 family farms (as defined by USDA), accounting for 97 percent of all farms and 89 percent of census farm area in the United States. In 1988 Mark Friedberger warned, "The farm family is a unique institution, perhaps the last remnant, in an increasingly complex world, of a simpler social order in which economic and domestic activities were inextricably bound together. In the past few years, however, American agriculture has suffered huge losses, and family farmers have seen their way of life threatened by economic forces beyond their control." However, by 1981 Ingolf Vogeler argued it was too late—the American family farm had been replaced by large agribusiness corporations pretending to be family operated.
A USDA survey conducted in 2011 estimated that family farms account for 85 percent of US farm production and 85 percent of US gross farm income. Mid-size and larger family farms account for 60 percent of US farm production and dominate US production of cotton, cash grain and hogs. Small family farms account for 26 percent of US farm production overall, and higher percentages of production of poultry, beef cattle, some other livestock and hay.
Several kinds of US family farms are recognized in USDA farm typology:
Small family farms are defined as those with annual gross cash farm income (GCFI) of less than $350,000; in 2011, these accounted for 90 percent of all US farms. Because low net farm incomes tend to predominate on such farms, most farm families on small family farms are extremely dependent on off-farm income.
Small family farms in which the principal operator was mostly employed off-farm accounted for 42 percent of all farms and 15 percent of total US farm area; median net farm income was $788.
Retirement family farms were small farms accounting for 16 percent of all farms and 7 percent of total US farm area; median net farm income was $5,002.
The other small family farm categories are those in which farming occupies at least 50 percent of the principal operator's working time. These are:
Low-sales small family farms (with GCFI less than $150,000); 26 percent of all US farms, 18 percent of total US farm area, median net farm income $3,579.
Moderate-sales small family farms (with GCFI of $150,000 to $349,999); 5.44 percent of all US farms, 13 percent of total US farm area, median net farm income $67,986.
Mid-size family farms (GCFI of $350,000 to $999,999); 6 percent of all US farms, 22 percent of total US farm area; median net farm income $154,538.
Large family farms (GCFI $1,000,000 to $4,999,999); 2 percent of all US farms, 14 percent of total US farm area; median net farm income $476,234.
Very large family farms (GCFI over $5,000,000); <1 percent of all US farms, 2 percent of total US farm area; median net farm income $1,910,454.
Family farms include not only sole proprietorships and family partnerships, but also family corporations. Family-owned corporations account for 5 percent of all farms and 89 percent of corporate farms in the United States. About 98 percent of US family corporations owning farms are small, with no more than 10 shareholders; average net farm income of family corporate farms was $189,400 in 2012. (In contrast, 90 percent of US non-family corporations owning farms are small, having no more than 10 shareholders; average net cash farm income for US non-family corporate farms was $270,670 in 2012.)
Canada
In Canada, the number of "family farms" cannot be inferred closely, because of the nature of census data, which do not distinguish family and non-family farm partnerships. In 2011, of Canada's 205,730 farms, 55 percent were sole proprietorships, 25 percent were partnerships, 17 percent were family corporations, 2 percent were non-family corporations and <1 percent were other categories. Because some but not all partnerships involve family members, these data suggest that family farms account for between about 73 and 97 percent of Canadian farms. The family farm percentage is likely to be near the high end of this range, for two reasons. The partners in a [Canadian] farm partnership are typically spouses, often forming the farm partnership for tax reasons. Also, as in the US, family farm succession planning can use a partnership as a means of apportioning family farm tenure among family members when a sole proprietor is ready to transfer some or all of ownership and operation of a farm to offspring. Conversion of a sole proprietorship family farm to a family corporation may also be influenced by legal and financial, e.g. tax, considerations. The Canadian Encyclopedia estimates that more than 90 percent of Canadian farms are family operations. In 2006, of Canadian farms with more than one million dollars in annual gross farm receipts, about 63 percent were family corporations and 13 percent were non-family corporations.
Europe
Analysis of data for 59,000 farms in the 12 member states of the European Community found that in 1989, about three-quarters of the farms were family farms, producing just over half of total agricultural output.
As of 2010, there were approximately 139,900 family farms in Ireland, with an average size of 35.7 hectares per holding. (Nearly all farms in Ireland are family farms.) In Ireland, average family farm income was 25,483 euros in 2012. Analysis by Teagasc (Ireland's Agriculture and Food Development Authority) estimates that 37 percent of Irish farms are economically viable and an additional 30 percent are sustainable due to income from off-farm sources; 33 percent meet neither criterion and are considered economically vulnerable.
Newly industrialized countries
In Brazil, there are about 4.37 million family farms. These account for 84.4 percent of farms, 24.3 percent of farmland area and 37.5 percent of the value of agricultural production.
Developing countries
In sub-Saharan Africa, 80% of the farms were family owned and worked by 2014.
Sub-Saharan agriculture was mostly defined by slash-and-burn subsistence farming, historically spread by the Bantu expansion.
Permanent farming estates were established during colonialism, in the 19th to 20th century.
After decolonisation, white farmers in some African countries have tended to be attacked, killed or evicted, notably in South Africa and Zimbabwe.
In southern Africa, "On peasant family farms ..., cash input costs are very low, non‐household labour is sourced largely from communal work groups through kinship ties, and support services needed to sustain production are minimal." On commercial family farms, "cash input costs are high, little non‐family labour is used and strong support services are necessary."
International Year of Family Farming
At the 66th session of the United Nations General Assembly, 2014 was formally declared to be the "International Year of Family Farming" (IYFF). The Food and Agriculture Organization of the United Nations was invited to facilitate its implementation, in collaboration with Governments, International Development Agencies, farmers' organizations and other relevant organizations of the United Nations system as well as relevant non-governmental organizations.
The goal of the 2014 IYFF is to reposition family farming at the centre of agricultural, environmental and social policies in the national agendas by identifying gaps and opportunities to promote a shift towards a more equal and balanced development. The 2014 IYFF will promote broad discussion and cooperation at the national, regional and global levels to increase awareness and understanding of the challenges faced by smallholders and help identify efficient ways to support family farmers.
| Technology | Agriculture, labor and economy | null |
171899 | https://en.wikipedia.org/wiki/Agathis%20australis | Agathis australis | Agathis australis, or kauri, is a coniferous tree in the family Araucariaceae, found north of 38°S in the northern regions of New Zealand's North Island.
It is the largest (by volume) but not tallest species of tree in New Zealand, standing up to 50 m tall in the emergent layer above the forest's main canopy. The tree has smooth bark and small narrow leaves. Other common names to distinguish A. australis from other members of Agathis are southern kauri and New Zealand kauri.
With its podsolization capability and regeneration pattern it can compete with faster growing angiosperms. Because it is such a conspicuous species, forest containing kauri is generally known as kauri forest, although kauri need not be the most abundant tree. In the warmer northern climate, kauri forests have a higher species richness than those found further south. Kauri even act as a foundation species that modify the soil under their canopy to create unique plant communities.
Taxonomy
Scottish botanist David Don described the species as Dammara australis. Agathis is derived from Greek and means 'ball of twine', a reference to the shape of the male cones, which are also known by the botanical term strobili. Australis translates in English to 'southern'.
Etymology
The Māori name is descended from Proto-Polynesian *kauquli, Samoan ebony or Diospyros samoensis.
Description
The young plant grows straight upwards and has the form of a narrow cone with branches going out along the length of the trunk. However, as it gains in height, the lowest branches are shed, preventing vines from climbing. By maturity, the top branches form an imposing crown that stands out over all other native trees, dominating the forest canopy.
The flaking bark of the kauri tree defends it from parasitic plants, and accumulates around the base of the trunk. On large trees it may pile up to a height of 2 m or more. The kauri has a habit of forming small clumps or patches scattered through mixed forests.
Kauri leaves are 3 to 7 cm long and 1 cm broad, tough and leathery in texture, with no midrib; they are arranged in opposite pairs or whorls of three on the stem. The seed cones are globose, 5 to 7 cm diameter, and mature 18 to 20 months after pollination; the seed cones disintegrate at maturity to release winged seeds, which are then dispersed by the wind. A single tree produces both male and female seed cones. Fertilisation of the seeds occurs by pollination, which may be driven by the same or another tree's pollen.
Size
Agathis australis can attain heights of 40 to 50 metres and trunk diameters big enough to rival Californian sequoias at over 5 metres. The largest kauri trees did not attain as much height or girth at ground level but contain more timber in their cylindrical trunks than comparable Sequoias with their tapering stems.
The largest recorded specimen was known as The Great Ghost and grew in the mountains at the head of the Tararu Creek, which drains into the Hauraki Gulf just north of the mouth of the Waihou River (Thames). Thames Historian Alastair Isdale says the tree was 8.54 metres in diameter, and 26.83 metres in girth. It was consumed by fire in .
A kauri tree at Mill Creek, Mercury Bay, known as Father of the Forests was measured in the early 1840s as 22 metres in circumference and 24 metres to the first branches. It was recorded as being killed by lightning in that period.
Another huge tree, Kairaru, had a girth of 20.1 metres and a columnar trunk free of branches for 30.5 metres as measured by a Crown Lands ranger, Henry Wilson, in 1860. It was on a spur of Mt Tutamoe about 30 km south of Waipoua Forest near Kaihau. It was destroyed in the 1880s or 1890s when a series of huge fires swept the area.
Other trees far larger than living kauri have been noted in other areas. Rumors of stumps up to 6 metres are sometimes suggested in areas such as the Billygoat Track above the Kauaeranga Valley near Thames. However, there is no good evidence for these (e.g., a documented measurement or a photograph with a person for scale).
Given that over 90 per cent of the area of kauri forest standing before 1000AD was destroyed by about 1900, it is not surprising that recent records are of smaller, but still very large trees. Two large kauri fell during tropical storms in the 1970s. One of these was Toronui, in Waipoua Forest. Its diameter was larger than that of Tāne Mahuta and its clean bole larger than that of Te Matua Ngahere, and by forestry measurements was the largest standing. Another tree, Kopi, in Omahuta Forest near the standing Hokianga kauri, was the third largest with a height of 56.39 metres (185') and a diameter of 4.19 metres (13.75'). It fell in 1973. Like many ancient kauri both trees were partly hollow.
Growth rate and age
In general over the lifetime of the tree the growth rate tends to increase, reach a maximum, then decline. A 1987 study measured mean annual diameter increments ranging from 1.5 to 4.6 mm per year with an overall average of 2.3 mm per year. This is equivalent to 8.7 annual rings per centimetre of core, said to be half the commonly quoted figure for growth rate. The same study found only a weak relationship between age and diameter. The growth of kauri in planted and second-growth natural forests has been reviewed and compared during the development of growth and yield models for the species. Kauri in planted forests were found to have up to 12 times the volume productivity than those in natural stands at the same age.
Individuals in the same 10 cm diameter class may vary in age by 300 years, and the largest individual on any particular site is often not the oldest. Trees can normally live longer than 600 years. Many individuals probably exceed 1000 years, but there is no conclusive evidence that trees can exceed 2000 years in age. By combining tree ring samples from living kauri, wooden buildings, and preserved swamp wood, a dendrochronology has been created which reaches back 4,500 years, the longest tree ring record of past climate change in the southern hemisphere. One 1700 year old swamp wood kauri that dates to approximately 42,000 years ago contains fine-scale carbon-14 fluctuations in its rings that may be reflective of the most recent magnetic field flip of the earth.
Root structure and soil interaction
Much like podocarps, it feeds in the organic litter near the surface of the soil through fine root hairs. This layer of the soil is composed of organic matter derived from falling leaves and branches as well as dead trees, and is constantly undergoing decomposition. On the other hand, broadleaf trees such as māhoe derive a good fraction of their nutrition in the deeper mineral layer of the soil. Although its feeding root system is very shallow, it also has several downwardly directed peg roots which anchor it firmly in the soil. Such a solid foundation is necessary to prevent a tree the size of a kauri from blowing over in storms and cyclones.
The litter left by kauri is much more acidic than most trees, and as it decays similarly acidic compounds are liberated. In a process known as leaching, these acidic molecules pass through the soil layers with the help of rainfall, and release other nutrients trapped in clay such as nitrogen and phosphorus. This leaves these important nutrients unavailable to other trees, as they are washed down into deeper layers. This process is known as podsolization, and changes the soil colour to a dull grey. For a single tree, this leaves an area of leached soil beneath known as a cup podsol (de). This leaching process is important for kauri's survival as it competes with other species for space.
Leaf litter and other decaying parts of a kauri decompose much more slowly than those of most other species. Besides its acidity, the plant also bears substances such as waxes and phenols, most notably tannins, that are harmful to microorganisms. This results in a large buildup of litter around the base of a mature tree in which its own roots feed. As with most perennials, these feeding roots also house a symbiotic fungi known as mycorrhiza which increase the plant's efficiency in taking up nutrients. In this mutualistic relationship, the fungus derives its own nutrition from the roots. In its interactions with the soil, kauri is thus able to starve its competitors of much needed nutrients and compete with much younger lineages.
The fungi on kauri are a food source for the larvae of the New Zealand giraffe weevil, Lasiorhynchus barbicornis. The larvae of L. barbicornis burrow into the wood of a tree for up to two years. Then L. barbicornis exit the bark of the tree as a fully formed adult beetle. These adult L. barbicornis exit from trees in Spring and Summer and months. After emerging from the tree, these adult L. barbicornis only live for a few weeks.
Distribution
Local spatial distribution
In terms of local topography, kauri is far from randomly dispersed. As mentioned above, kauri relies on depriving its competitors of nutrition in order to survive. However, one important consideration not discussed thus far is the slope of the land. Water on hills flows downward by the action of gravity, taking with it the nutrients in the soil. This results in a gradient from nutrient poor soil at the top of slopes to nutrient rich soils below. As nutrients leached are replaced by aqueous nitrates and phosphates from above, the kauri tree is less able to inhibit the growth of strong competitors such as angiosperms. In contrast, the leaching process is only enhanced on higher elevation. In Waipoua Forest this is reflected in higher abundances of kauri on ridge crests, and greater concentrations of its main competitors, such as tarairi, at low elevations. This pattern is known as niche partitioning, and allows more than one species to occupy the same area. Those species which live alongside kauri include tawari, a montane broadleaf tree which is normally found in higher altitudes, where nutrient cycling is naturally slow.
Changes over recent geological time
Kauri is found growing in its natural ecosystem north of 38°S latitude. Its southern limit stretches from the Kawhia Harbour in the west to the eastern Kaimai Range.
However, its distribution has changed greatly over geological time because of climate change. This is shown in the recent Holocene epoch by its migration southwards after the peak of the last ice age. During this time when frozen ice sheets covered much of the world's continents, kauri was able to survive only in isolated pockets, its main refuge being in the very far north. Radiocarbon dating is one technique used by scientists to uncover the history of the tree's distribution, with stump kauri from peat swamps used for measurement. The coldest period in recent times occurred about 15,000 to 20,000 years ago, during which time kauri was apparently confined north of Kaitaia, near the northernmost point of the North Island, North Cape. Kauri requires a mean temperature of 17 °C or more for most of the year. The tree's retreat can be used as a proxy for temperature changes during this period. While not present in modern days, the Aupōuri Peninsula in the far north was a refuge for kauri, as large quantities of kauri gum were present in the soils.
It remains unclear whether kauri recolonised the North Island from a single refuge in the far north or from scattered pockets of isolated stands that managed to survive despite climatic conditions. It spread south through Whangārei, past Dargaville and as far south as Waikato, attaining its peak distribution during the years 3000 BP to 2000 BP. There is some suggestion that it has receded somewhat since then, which may indicate temperatures have declined slightly. During the peak of its movement southwards, it was travelling as fast as 200 metres per year. Its southward spread seems relatively rapid for a tree that can take a millennium to reach complete maturity. This can be explained by its life history pattern.
Kauri relies on wind for pollination and seed dispersal, while many other native trees have their seeds carried large distances by frugivores (animals which eat fruit) such as the kererū (native pigeon). However, kauri trees can produce seeds while relatively young, taking only 50 years or so before giving rise to their own offspring. This trait makes them somewhat like a pioneer species, despite the fact that their long lifespan is characteristic of K-selected species. In good conditions, where access to water and sunlight are above average, diameters in excess of 15 centimetres and seed production can occur inside 15 years.
Regeneration and life history
Just as the niche of kauri is differentiated through its interactions with the soil, it also has a separate regeneration 'strategy' compared to its broadleaf neighbours. The relationship is very similar to the podocarp-broadleaf forests further south. Kauri demand much more light and require larger gaps to regenerate than such broadleaf trees as pūriri and kohekohe that show far more shade tolerance. Unlike kauri, these broadleaf species can regenerate in areas where lower levels of light reach ground level, for example from a single branch falling off. Kauri trees must therefore remain alive long enough for a large disturbance to occur, allowing them sufficient light to regenerate. In areas where large amounts of forest are destroyed, such as by logging, kauri seedlings are able to regenerate much more easily due not only to increased sunlight, but their relatively strong resistance to wind and frosts. Kauri occupy the emergent layer of the forest, where they are exposed to the effects of the weather; however, the smaller trees that dominate the main canopy are sheltered both by the emergent trees above and by each other. Left in open areas without protection, these smaller trees are far less capable of regenerating.
When there is a disturbance severe enough to favour their regeneration, kauri trees regenerate en masse, producing a generation of trees of similar age after each disturbance. The distribution of kauri allows researchers to deduce when and where disturbances have occurred, and how large they may have been; the presence of abundant kauri may indicate that an area is prone to disturbance. Kauri seedlings can still occur in areas with low light but mortality rates increase for such seedlings, and those that survive self-thinning and grow to sapling stage tend to be found in higher light environments.
During periods with less disturbance kauri tends to lose ground to broadleaf competitors which can better tolerate shaded environments. In the complete absence of disturbance, kauri tends to become rare as it is excluded by its competitors. Kauri biomass tends to decrease during such times, as more biomass becomes concentrated in angiosperm species like tōwai. Kauri trees also tend to become more randomly distributed in age, with each tree dying at a different point in time, and regeneration gaps becoming rare and sporadic. Over thousands of years these varying regeneration strategies produce a tug of war effect where kauri retreats uphill during periods of calm, then takes over lower areas briefly during mass disturbances. Although such trends cannot be observed in a human lifetime, research into current patterns of distribution, behaviour of species in experimental conditions, and study of pollen sediments (see palynology) have helped shed light on the life history of kauri.
Kauri seeds may generally be taken from mature cones in late March. Each scale on a cone contains a single winged seed approximately 5 mm by 8 mm and attached to a thin wing perhaps half as large again. The cone is fully open and dispersed within only two to three days of starting.
Studies show that kauri develop root grafts through which they share water and nutrients with neighbours of the same species.
Ethnobotany
Deforestation
Heavy logging, which began around 1820 and continued for a century, has considerably decreased the number of kauri trees. It has been estimated that before 1840, the kauri forests of northern New Zealand occupied at least 12,000 square kilometres. The British Royal Navy sent four vessels, HMS Coromandel (1821), HMS Dromedary (1821), (1840), and HMS Tortoise (1841) to gather kauri-wood spars.
By 1900, less than 10 per cent of the original kauri survived. By the 1950s this area had decreased to about 1,400 square kilometres in 47 forests depleted of their best kauri. It is estimated that today, there is 4 per cent of uncut forest left in small pockets.
Estimates are that around half of the timber was accidentally or deliberately burnt. More than half of the remainder had been exported to Australia, Britain, and other countries, while the balance was used locally to build houses and ships. Much of the timber was sold for a return sufficient only to cover wages and expenses. From 1871 to 1895 the receipts indicate a rate of about 8 shillings (around NZ$20 in 2003) per 100 superficial feet (34 shillings/m3).
The Government continued to sell large areas of kauri forests to sawmillers who, under no restrictions, took the most effective and economical steps to secure the timber, resulting in much waste and destruction. At a sale in 1908 more than 5,000 standing kauri trees, totalling about 20,000,000 superficial feet (47,000 m3), were sold for less than £2 per tree (£2 in 1908 equates to around NZ$100 in 2003). It is said that in 1890 the royalty on standing timber fell in some cases to as low as twopence (NZ$0.45 in 2003) per 100 superficial feet (8 pence/m3), though the expense of cutting and removing it to the mills was typically great due to the difficult terrain where they were located.
Probably the most controversial kauri logging decision in the last century was that of the National Government to initiate clear fell logging of the Warawara state forest (North of the Hokianga) in the late 1960s. This created a national outcry as this forest contains the second largest volume of kauri after the Waipoua forest and was until that time, essentially unlogged (Adams, 1980). The plan also involved considerable cost, requiring a long road to be driven up a steep high plateau into the heart of the protected area. Because the stands of kauri were dense, the ecological destruction in the affected plateau area (approximately a fifth of the forest by area, and a quarter by volume of timber) was essentially complete (as of the early 1990s most of the affected area contained a thick covering of native grasses with little or no kauri regeneration). Logging was stopped in fulfillment of an election pledge by the Labour Government of 1972. When the National Party was reelected in 1975, the ban on kauri logging in the Warawara remained in place, but was soon replaced by policies encouraging the logging of giant tōtara and other podocarps in the central North Island. The outcry over the Warawara was an important stepping stone towards the legal protection of the small percentage of remaining virgin kauri-podocarp forest in New Zealand's Government-owned forests.
Uses
Although today its use is far more restricted, in the past the size and strength of kauri timber made it a popular wood for construction and ship building, particularly for masts of sailing ships because of its parallel grain and the absence of branches for much of its height. Kauri crown and stump wood was much appreciated for its beauty, and was sought after for ornamental wood panelling as well as high-end furniture. Although not as highly prized, the light colour of kauri trunk wood made it also well-suited for more utilitarian furniture construction, as well as for use in the fabrication of cisterns, barrels, bridge construction material, fences, moulds for metal forges, large rollers for the textile industry, railway sleepers and cross bracing for mines and tunnels.
In the late 19th and early 20th centuries Kauri gum (semi-fossilised kauri resin) was a valuable commodity, particularly for varnish, spurring the development of a gum-digger industry.
Today, the kauri is being considered as a long-term carbon sink. This is because estimates of the total carbon content in living above ground biomass and dead biomass of mature kauri forest are the second highest of any forest type recorded anywhere in the world. The estimated total carbon capture is up to nearly 1000 tonnes per hectare. In this capacity, kauri are bettered only by mature Eucalyptus regnans forest, and are far higher than any tropical or boreal forest type yet recorded. It is also conjectured that the process of carbon capture does not reach equilibrium, which along with no need of direct maintenance, makes kauri forests a potentially attractive alternative to short rotation forestry options such as Pinus radiata.
Timber
Technical specifications
Moisture content of dried wood: 12 per cent
Density of wood: 560 kg/m3
Tensile strength: 88 MPa
Modulus of elasticity: 9.1 GPa
After felled kauri wood dries to a 12 per cent moisture content, the tangential contraction is 4.1 per cent and the radial contraction is 2.3 per cent.
Kauri is considered a first rate timber. The whiter sapwood is generally slightly lighter in weight. Kauri is not highly resistant to rot and when used in boatbuilding must be protected from the elements with paint, varnish or epoxy to avoid rot. Its popularity with boatbuilders is due to its very long, clear lengths, its relatively light weight and its beautiful sheen when oiled or varnished. Kauri wood planes and saws easily. Its wood holds screws and nails very well and does not readily split, crack, or warp. Kauri wood darkens with age to a richer golden brown colour. Very little New Zealand kauri is now sold, and the most commonly available kauri in New Zealand is Fiji kauri, which is very similar in appearance but lighter in weight.
Swamp kauri
Prehistoric kauri forests have been preserved in waterlogged soils as swamp kauri. A considerable number of kauri have been found buried in salt marshes, resulting from ancient natural changes such as volcanic eruptions, sea-level changes and floods. Such trees have been radiocarbon dated to 50,000 years ago or older. The bark and the seed cones of the trees often survive together with the trunk, although when excavated and exposed to the air, these parts undergo rapid deterioration. The quality of the disinterred wood varies. Some is in good shape, comparable to that of newly felled kauri, although often lighter in colour. The colour can be improved by the use of natural wood stains to heighten the details of the grain. After a drying process, such ancient kauri can be used for furniture, but not for construction.
Conservation
The small remaining pockets of kauri forest in New Zealand have survived in areas that were not subjected to burning by Māori and were too inaccessible for European loggers. The largest area of mature kauri forest is Waipoua Forest in Northland. Mature and regenerating kauri can also be found in other National and Regional Parks such as Puketi and Omahuta Forests in Northland, the Waitākere Ranges near Auckland, and Coromandel Forest Park on the Coromandel Peninsula.
The importance of Waipoua Forest in relation to the kauri was that it remained the only kauri forest retaining its former virgin condition, and that it was extensive enough to give reasonable promise of permanent survival. On 2 July 1952 an area of over 80 km2 of Waipoua was proclaimed a forest sanctuary after a petition to the Government.
The zoologist William Roy McGregor was one of the driving forces in this movement, writing an 80-page illustrated pamphlet on the subject, which proved an effective manifesto for conservation.
Along with the Warawara to the North, Waipoua Forest contains three quarters of New Zealand's remaining kauri. Kauri Grove on the Coromandel Peninsula is another area with a remaining cluster of kauri, and includes the Siamese Kauri, two trees with a conjoined lower trunk.
In 1921 a philanthropic Cornishman named James Trounson sold to the Government for £40,000, a large area adjacent to a few acres of Crown land and said to contain at least 4,000 kauri trees. From time to time Trounson gifted additional land, until what is known as Trounson Park comprised a total of 4 km2.
The most famous specimens are Tāne Mahuta and Te Matua Ngahere in Waipoua Forest. These two trees have become tourist attractions because of their size and accessibility. Tane Mahuta, named after the Māori forest god, is the biggest existing kauri with a girth of , a trunk height of , a total height of and a total volume including the crown of . Te Matua Ngahere, which means 'Father of the Forest', is smaller but stouter than Tane Mahuta, with a girth (circumference) of . Important note: all the measurements above were taken in 1971.
Kauri is common as a specimen tree in parks and gardens throughout New Zealand, prized for the distinctive look of young trees, its low maintenance once established (although seedlings are frost tender).
Kauri dieback
Kauri dieback was observed in the Waitākere Ranges caused by Phytophthora cinnamomi in the 1950s, again on Great Barrier Island in 1972 linked to a different pathogen, Phytophthora agathidicida and subsequently spread to kauri forest on the mainland. The disease, known as kauri dieback or kauri collar rot, is believed to be over 300 years old and causes yellowing leaves, thinning canopy, dead branches, lesions that bleed resin, and tree death.
Phytophthora agathidicida was identified as a new species in April 2008. Its closest known relative is Phytophthora katsurae. The pathogen is believed to be spread on people's shoes or by mammals, particularly feral pigs. A collaborative response team has been formed to work on the disease. The team includes MAF Biosecurity, the Conservation Department, Auckland and Northland regional councils, Waikato Regional Council, and Bay of Plenty Regional Council. The team is charged with assessing the risk, determining methods and their feasibility to limit the spread, collecting more information (e.g. how widespread), and ensuring a coordinated response. The Department of Conservation has issued guidelines to prevent the spread of the disease, including keeping to defined tracks, cleaning footwear before and after entering kauri forest areas, and staying away from kauri roots.
| Biology and health sciences | Pinophyta (Conifers) | Plants |
171950 | https://en.wikipedia.org/wiki/Root%20of%20unity | Root of unity | In mathematics, a root of unity, occasionally called a de Moivre number, is any complex number that yields 1 when raised to some positive integer power . Roots of unity are used in many branches of mathematics, and are especially important in number theory, the theory of group characters, and the discrete Fourier transform.
Roots of unity can be defined in any field. If the characteristic of the field is zero, the roots are complex numbers that are also algebraic integers. For fields with a positive characteristic, the roots belong to a finite field, and, conversely, every nonzero element of a finite field is a root of unity. Any algebraically closed field contains exactly th roots of unity, except when is a multiple of the (positive) characteristic of the field.
General definition
An th root of unity, where is a positive integer, is a number satisfying the equation
Unless otherwise specified, the roots of unity may be taken to be complex numbers (including the number 1, and the number −1 if is even, which are complex with a zero imaginary part), and in this case, the th roots of unity are
However, the defining equation of roots of unity is meaningful over any field (and even over any ring) , and this allows considering roots of unity in . Whichever is the field , the roots of unity in are either complex numbers, if the characteristic of is 0, or, otherwise, belong to a finite field. Conversely, every nonzero element in a finite field is a root of unity in that field. See Root of unity modulo n and Finite field for further details.
An th root of unity is said to be if it is not an th root of unity for some smaller , that is if
If n is a prime number, then all th roots of unity, except 1, are primitive.
In the above formula in terms of exponential and trigonometric functions, the primitive th roots of unity are those for which and are coprime integers.
Subsequent sections of this article will comply with complex roots of unity. For the case of roots of unity in fields of nonzero characteristic, see . For the case of roots of unity in rings of modular integers, see Root of unity modulo n.
Elementary properties
Every th root of unity is a primitive th root of unity for some , which is the smallest positive integer such that .
Any integer power of an th root of unity is also an th root of unity, as
This is also true for negative exponents. In particular, the reciprocal of an th root of unity is its complex conjugate, and is also an th root of unity:
If is an th root of unity and then . Indeed, by the definition of congruence modulo n, for some integer , and hence
Therefore, given a power of , one has , where is the remainder of the Euclidean division of by .
Let be a primitive th root of unity. Then the powers , , ..., , are th roots of unity and are all distinct. (If where , then , which would imply that would not be primitive.) This implies that , , ..., , are all of the th roots of unity, since an th-degree polynomial equation over a field (in this case the field of complex numbers) has at most solutions.
From the preceding, it follows that, if is a primitive th root of unity, then if and only if
If is not primitive then implies but the converse may be false, as shown by the following example. If , a non-primitive th root of unity is , and one has , although
Let be a primitive th root of unity. A power of is a primitive th root of unity for
where is the greatest common divisor of and . This results from the fact that is the smallest multiple of that is also a multiple of . In other words, is the least common multiple of and . Thus
Thus, if and are coprime, is also a primitive th root of unity, and therefore there are distinct primitive th roots of unity (where is Euler's totient function). This implies that if is a prime number, all the roots except are primitive.
In other words, if is the set of all th roots of unity and is the set of primitive ones, is a disjoint union of the :
where the notation means that goes through all the positive divisors of , including and .
Since the cardinality of is , and that of is , this demonstrates the classical formula
Group properties
Group of all roots of unity
The product and the multiplicative inverse of two roots of unity are also roots of unity. In fact, if and , then , and , where is the least common multiple of and .
Therefore, the roots of unity form an abelian group under multiplication. This group is the torsion subgroup of the circle group.
Group of th roots of unity
For an integer n, the product and the multiplicative inverse of two th roots of unity are also th roots of unity. Therefore, the th roots of unity form an abelian group under multiplication.
Given a primitive th root of unity , the other th roots are powers of . This means that the group of the th roots of unity is a cyclic group. It is worth remarking that the term of cyclic group originated from the fact that this group is a subgroup of the circle group.
Galois group of the primitive th roots of unity
Let be the field extension of the rational numbers generated over by a primitive th root of unity . As every th root of unity is a power of , the field contains all th roots of unity, and is a Galois extension of
If is an integer, is a primitive th root of unity if and only if and are coprime. In this case, the map
induces an automorphism of , which maps every th root of unity to its th power. Every automorphism of is obtained in this way, and these automorphisms form the Galois group of over the field of the rationals.
The rules of exponentiation imply that the composition of two such automorphisms is obtained by multiplying the exponents. It follows that the map
defines a group isomorphism between the units of the ring of integers modulo and the Galois group of
This shows that this Galois group is abelian, and implies thus that the primitive roots of unity may be expressed in terms of radicals.
Galois group of the real part of the primitive roots of unity
The real part of the primitive roots of unity are related to one another as roots of the minimal polynomial of The roots of the minimal polynomial are just twice the real part; these roots form a cyclic Galois group.
Trigonometric expression
De Moivre's formula, which is valid for all real and integers , is
Setting gives a primitive th root of unity – one gets
but
for . In other words,
is a primitive th root of unity.
This formula shows that in the complex plane the th roots of unity are at the vertices of a regular -sided polygon inscribed in the unit circle, with one vertex at 1 (see the plots for and on the right). This geometric fact accounts for the term "cyclotomic" in such phrases as cyclotomic field and cyclotomic polynomial; it is from the Greek roots "cyclo" (circle) plus "tomos" (cut, divide).
Euler's formula
which is valid for all real , can be used to put the formula for the th roots of unity into the form
It follows from the discussion in the previous section that this is a primitive th-root if and only if the fraction is in lowest terms; that is, that and are coprime. An irrational number that can be expressed as the real part of the root of unity; that is, as , is called a trigonometric number.
Algebraic expression
The th roots of unity are, by definition, the roots of the polynomial , and are thus algebraic numbers. As this polynomial is not irreducible (except for ), the primitive th roots of unity are roots of an irreducible polynomial (over the integers) of lower degree, called the th cyclotomic polynomial, and often denoted . The degree of is given by Euler's totient function, which counts (among other things) the number of primitive th roots of unity. The roots of are exactly the primitive th roots of unity.
Galois theory can be used to show that the cyclotomic polynomials may be conveniently solved in terms of radicals. (The trivial form is not convenient, because it contains non-primitive roots, such as 1, which are not roots of the cyclotomic polynomial, and because it does not give the real and imaginary parts separately.) This means that, for each positive integer , there exists an expression built from integers by root extractions, additions, subtractions, multiplications, and divisions (and nothing else), such that the primitive th roots of unity are exactly the set of values that can be obtained by choosing values for the root extractions ( possible values for a th root). (For more details see , below.)
Gauss proved that a primitive th root of unity can be expressed using only square roots, addition, subtraction, multiplication and division if and only if it is possible to construct with compass and straightedge the regular -gon. This is the case if and only if is either a power of two or the product of a power of two and Fermat primes that are all different.
If is a primitive th root of unity, the same is true for , and is twice the real part of . In other words, is a reciprocal polynomial, the polynomial that has as a root may be deduced from by the standard manipulation on reciprocal polynomials, and the primitive th roots of unity may be deduced from the roots of by solving the quadratic equation That is, the real part of the primitive root is and its imaginary part is
The polynomial is an irreducible polynomial whose roots are all real. Its degree is a power of two, if and only if is a product of a power of two by a product (possibly empty) of distinct Fermat primes, and the regular -gon is constructible with compass and straightedge. Otherwise, it is solvable in radicals, but one are in the casus irreducibilis, that is, every expression of the roots in terms of radicals involves nonreal radicals.
Explicit expressions in low degrees
For , the cyclotomic polynomial is Therefore, the only primitive first root of unity is 1, which is a non-primitive th root of unity for every n > 1.
As , the only primitive second (square) root of unity is −1, which is also a non-primitive th root of unity for every even . With the preceding case, this completes the list of real roots of unity.
As , the primitive third (cube) roots of unity, which are the roots of this quadratic polynomial, are
As , the two primitive fourth roots of unity are and .
As , the four primitive fifth roots of unity are the roots of this quartic polynomial, which may be explicitly solved in terms of radicals, giving the roots where may take the two values 1 and −1 (the same value in the two occurrences).
As , there are two primitive sixth roots of unity, which are the negatives (and also the square roots) of the two primitive cube roots:
As 7 is not a Fermat prime, the seventh roots of unity are the first that require cube roots. There are 6 primitive seventh roots of unity, which are pairwise complex conjugate. The sum of a root and its conjugate is twice its real part. These three sums are the three real roots of the cubic polynomial and the primitive seventh roots of unity are where runs over the roots of the above polynomial. As for every cubic polynomial, these roots may be expressed in terms of square and cube roots. However, as these three roots are all real, this is casus irreducibilis, and any such expression involves non-real cube roots.
As , the four primitive eighth roots of unity are the square roots of the primitive fourth roots, . They are thus
See Heptadecagon for the real part of a 17th root of unity.
Periodicity
If is a primitive th root of unity, then the sequence of powers
is -periodic (because for all values of ), and the sequences of powers
for are all -periodic (because ). Furthermore, the set } of these sequences is a basis of the linear space of all -periodic sequences. This means that any -periodic sequence of complex numbers
can be expressed as a linear combination of powers of a primitive th root of unity:
for some complex numbers and every integer .
This is a form of Fourier analysis. If is a (discrete) time variable, then is a frequency and is a complex amplitude.
Choosing for the primitive th root of unity
allows to be expressed as a linear combination of and :
This is a discrete Fourier transform.
Summation
Let be the sum of all the th roots of unity, primitive or not. Then
This is an immediate consequence of Vieta's formulas. In fact, the th roots of unity being the roots of the polynomial , their sum is the coefficient of degree , which is either 1 or 0 according whether or .
Alternatively, for there is nothing to prove, and for there exists a root – since the set of all the th roots of unity is a group, , so the sum satisfies , whence .
Let be the sum of all the primitive th roots of unity. Then
where is the Möbius function.
In the section Elementary properties, it was shown that if is the set of all th roots of unity and is the set of primitive ones, is a disjoint union of the :
This implies
Applying the Möbius inversion formula gives
In this formula, if , then , and for : . Therefore, .
This is the special case of Ramanujan's sum , defined as the sum of the th powers of the primitive th roots of unity:
Orthogonality
From the summation formula follows an orthogonality relationship: for and
where is the Kronecker delta and is any primitive th root of unity.
The matrix whose th entry is
defines a discrete Fourier transform. Computing the inverse transformation using Gaussian elimination requires operations. However, it follows from the orthogonality that is unitary. That is,
and thus the inverse of is simply the complex conjugate. (This fact was first noted by Gauss when solving the problem of trigonometric interpolation.) The straightforward application of or its inverse to a given vector requires operations. The fast Fourier transform algorithms reduces the number of operations further to .
Cyclotomic polynomials
The zeros of the polynomial
are precisely the th roots of unity, each with multiplicity 1. The th cyclotomic polynomial is defined by the fact that its zeros are precisely the primitive th roots of unity, each with multiplicity 1.
where are the primitive th roots of unity, and is Euler's totient function. The polynomial has integer coefficients and is an irreducible polynomial over the rational numbers (that is, it cannot be written as the product of two positive-degree polynomials with rational coefficients). The case of prime , which is easier than the general assertion, follows by applying Eisenstein's criterion to the polynomial
and expanding via the binomial theorem.
Every th root of unity is a primitive th root of unity for exactly one positive divisor of . This implies that
This formula represents the factorization of the polynomial into irreducible factors:
Applying Möbius inversion to the formula gives
where is the Möbius function. So the first few cyclotomic polynomials are
If is a prime number, then all the th roots of unity except 1 are primitive th roots. Therefore,
Substituting any positive integer ≥ 2 for , this sum becomes a base repunit. Thus a necessary (but not sufficient) condition for a repunit to be prime is that its length be prime.
Note that, contrary to first appearances, not all coefficients of all cyclotomic polynomials are 0, 1, or −1. The first exception is . It is not a surprise it takes this long to get an example, because the behavior of the coefficients depends not so much on as on how many odd prime factors appear in . More precisely, it can be shown that if has 1 or 2 odd prime factors (for example, ) then the th cyclotomic polynomial only has coefficients 0, 1 or −1. Thus the first conceivable for which there could be a coefficient besides 0, 1, or −1 is a product of the three smallest odd primes, and that is . This by itself doesn't prove the 105th polynomial has another coefficient, but does show it is the first one which even has a chance of working (and then a computation of the coefficients shows it does). A theorem of Schur says that there are cyclotomic polynomials with coefficients arbitrarily large in absolute value. In particular, if where are odd primes, and t is odd, then occurs as a coefficient in the th cyclotomic polynomial.
Many restrictions are known about the values that cyclotomic polynomials can assume at integer values. For example, if is prime, then if and only if .
Cyclotomic polynomials are solvable in radicals, as roots of unity are themselves radicals. Moreover, there exist more informative radical expressions for th roots of unity with the additional property that every value of the expression obtained by choosing values of the radicals (for example, signs of square roots) is a primitive th root of unity. This was already shown by Gauss in 1797. Efficient algorithms exist for calculating such expressions.
Cyclic groups
The th roots of unity form under multiplication a cyclic group of order , and in fact these groups comprise all of the finite subgroups of the multiplicative group of the complex number field. A generator for this cyclic group is a primitive th root of unity.
The th roots of unity form an irreducible representation of any cyclic group of order . The orthogonality relationship also follows from group-theoretic principles as described in Character group.
The roots of unity appear as entries of the eigenvectors of any circulant matrix; that is, matrices that are invariant under cyclic shifts, a fact that also follows from group representation theory as a variant of Bloch's theorem. In particular, if a circulant Hermitian matrix is considered (for example, a discretized one-dimensional Laplacian with periodic boundaries), the orthogonality property immediately follows from the usual orthogonality of eigenvectors of Hermitian matrices.
Cyclotomic fields
By adjoining a primitive th root of unity to one obtains the th cyclotomic field This field contains all th roots of unity and is the splitting field of the th cyclotomic polynomial over The field extension has degree φ(n) and its Galois group is naturally isomorphic to the multiplicative group of units of the ring
As the Galois group of is abelian, this is an abelian extension. Every subfield of a cyclotomic field is an abelian extension of the rationals. It follows that every nth root of unity may be expressed in term of k-roots, with various k not exceeding φ(n). In these cases Galois theory can be written out explicitly in terms of Gaussian periods: this theory from the Disquisitiones Arithmeticae of Gauss was published many years before Galois.
Conversely, every abelian extension of the rationals is such a subfield of a cyclotomic field – this is the content of a theorem of Kronecker, usually called the Kronecker–Weber theorem on the grounds that Weber completed the proof.
Relation to quadratic integers
For , both roots of unity and are integers.
For three values of , the roots of unity are quadratic integers:
For they are Eisenstein integers ().
For they are Gaussian integers (): see Imaginary unit.
For four other values of , the primitive roots of unity are not quadratic integers, but the sum of any root of unity with its complex conjugate (also an th root of unity) is a quadratic integer.
For , none of the non-real roots of unity (which satisfy a quartic equation) is a quadratic integer, but the sum of each root with its complex conjugate (also a 5th root of unity) is an element of the ring (). For two pairs of non-real 5th roots of unity these sums are inverse golden ratio and minus golden ratio.
For , for any root of unity equals to either 0, ±2, or ± ().
For , for any root of unity, equals to either 0, ±1, ±2 or ± ().
| Mathematics | Complex analysis | null |
172081 | https://en.wikipedia.org/wiki/Gooseberry | Gooseberry | Gooseberry ( or (American and northern British) or (southern British)) is a common name for many species of Ribes (which also includes currants), as well as a large number of plants of similar appearance, and also several unrelated plants (see List of gooseberries). The berries of those in the genus Ribes (sometimes placed in the genus Grossularia) are edible and may be green, orange, red, purple, yellow, white, or black.
Etymology
The goose in gooseberry has been mistakenly seen as a corruption of either the Dutch word or the allied German , or of the earlier forms of the French . Alternatively, the word has been connected to the Middle High German ('curl, crisped'), in Latin as . However, the Oxford English Dictionary takes the more literal derivation from goose and berry as probable because "the grounds on which plants and fruits have received names associating them with animals are so often inexplicable that the inappropriateness in the meaning does not necessarily give good grounds for believing that the word is an etymological corruption". The French for gooseberry is , translated as 'mackerel berries', due to their use in a sauce for mackerel in old French cuisine. The word first appears in written English in the 16th Century. In Britain, gooseberries may informally be called goosegogs.
Gooseberry bush was 19th-century slang for pubic hair, and from this comes the saying that babies are "born under a gooseberry bush".
Ecology
Black bears, various birds and small mammals eat the berries, while game animals, coyotes, foxes and raccoons browse the foliage.
Cultivation
In history
Gooseberry growing was popular in 19th-century Britain. The 1879 edition of the Encyclopedia Britannica described gooseberries thus:
The gooseberry is indigenous to many parts of Europe and western Asia, growing naturally in alpine thickets and rocky woods in the lower country, from France eastward, well into the Himalayas and peninsular India.
In Britain, it is often found in copses and hedgerows and about old ruins, but the gooseberry has been cultivated for so long that it is difficult to distinguish wild bushes from feral ones, or to determine where the gooseberry fits into the native flora of the island. Common as it is now on some of the lower slopes of the Alps of Piedmont and Savoy, it is uncertain whether the Romans were acquainted with the gooseberry, though it may possibly be alluded to in a vague passage of Pliny the Elder's Natural History; the hot summers of Italy, in ancient times as at present, would be unfavourable to its cultivation. Although gooseberries are now abundant in Germany and France, it does not appear to have been much grown there in the Middle Ages, though the wild fruit was held in some esteem medicinally for the cooling properties of its acid juice in fevers; while the old English name, Fea-berry, still surviving in some provincial dialects, indicates that it was similarly valued in Britain, where it was planted in gardens at a comparatively early period.
William Turner describes the gooseberry in his Herball, written about the middle of the 16th century, and a few years later it is mentioned in one of Thomas Tusser's quaint rhymes as an ordinary object of garden culture. Improved varieties were probably first raised by the skilful gardeners of Holland, whose name for the fruit, Kruisbezie, may have been corrupted into the present English vernacular word. Towards the end of the 18th century the gooseberry became a favourite object of cottage-horticulture, especially in Lancashire, where the working cotton-spinners raised numerous varieties from seed, their efforts having been chiefly directed to increasing the size of the fruit.
Of the many hundred sorts enumerated in recent horticultural works, few perhaps equal in flavour some of the older denizens of the fruit-garden, such as the Old Rough Red and Hairy Amber. The climate of the British Isles seems peculiarly adapted to bring the gooseberry to perfection, and it may be grown successfully even in the most northern parts of Scotland; indeed, the flavour of the fruit is said to improve with increasing latitude. In Norway even, the bush flourishes in gardens on the west coast nearly up to the Arctic Circle, and it is found wild as far north as 63°. The dry summers of the French and German plains are less suited to it, though it is grown in some hilly districts with tolerable success. The gooseberry in the south of England will grow well in cool situations and may sometimes be seen in gardens near London flourishing under the partial shade of apple trees, but in the north it needs full exposure to the sun to bring the fruit to perfection. It will succeed in almost any soil but prefers a rich loam or black alluvium, and, though naturally a plant of rather dry places, will do well in moist land, if drained.
The gooseberry was more populous in North America before it was discovered that it carries blister rust, deadly to certain pines, resulting in its removal from forest areas.
Modern cultivation
Humans cultivate gooseberries as insect habitats or directly for the sweet fruits. Numerous cultivars have been developed for both commercial and domestic use. Of special note are Ribes 'Careless', 'Greenfinch', 'Invicta', 'Leveller', and 'Whinham's Industry', to which the Royal Horticultural Society has awarded Garden Merit.
Ribes gooseberries are commonly raised from cuttings rather than seed; cuttings planted in the autumn will take root quickly and begin to bear fruit within a few years. Nevertheless, bushes planted from seed also rapidly reach maturity, exhibit similar pest-tolerance, and yield heavily. Fruit is produced on lateral spurs and the previous year's shoots.
Gooseberries must be pruned to insolate the interior and make space for the next year's branches, as well as reduce scratching from the spines when picking. Overladen branches can be (and often are) cut off complete with berries without substantially harming the plant. Heavy nitrogen composting produces excessive growth, weakening the bush to mildew.
Fungal pests
Gooseberries, like other members of genus Ribes, are banned or restricted in several states of the United States because they are secondary (telial) hosts for white pine blister rust.
Insect habitat
Gooseberry bushes (Ribes) are hosts to magpie moth (Abraxas grossulariata) caterpillars. Gooseberry plants are also a preferred host plant for comma butterfly (Polygonia c-album), whose larvae frequently feed upon the plant during the development stage, v-moth (Macaria wauaria), and gooseberry sawfly (Nematus ribesii). Nematus ribesii grubs will bury themselves in the ground to pupate; on hatching into adult form, they lay their eggs, which hatch into larvae on the underside of gooseberry leaves.
Culinary uses
Gooseberries are edible and can be eaten raw, or cooked as an ingredient in desserts, such as pies, fools and crumbles. Early pickings are generally sour and more appropriate for culinary use. This includes most supermarket gooseberries, which are often picked before fully ripe to increase shelf life. Gooseberries are also used to flavour beverages such as sodas, flavoured waters, or milk, and can be made into fruit wines and teas. Gooseberries can be preserved in the form of jams, dried fruit, as the primary or a secondary ingredient in pickling, or stored in sugar syrup. Pastry dishes often pair gooseberry with flavors such as hazelnut, honey, raspberry, strawberry, and white chocolate.
Nutrition
Raw gooseberries are 88% water, 10% carbohydrates, and 1% each of fat and protein. In a reference amount of , raw gooseberries supply 44 calories and are a rich source of vitamin C (31% of the Daily Value), with no other micronutrients in significant content.
| Biology and health sciences | Berries | Plants |
172111 | https://en.wikipedia.org/wiki/Washing%20machine | Washing machine | A washing machine (laundry machine, clothes washer, washer, or simply wash) is a machine designed to launder clothing. The term is mostly applied to machines that use water. Other ways of doing laundry include dry cleaning (which uses alternative cleaning fluids and is performed by specialist businesses) and ultrasonic cleaning.
Modern-day home appliances use electric power to automatically clean clothes. The user adds laundry detergent, which is sold in liquid, powder, or dehydrated sheet form, to the wash water. The machines are also found in commercial laundromats where customers pay-per-use.
History
Washing by hand
Laundering by hand involves soaking, beating, scrubbing, and rinsing dirty textiles. Before indoor plumbing, it was necessary to carry all the water used for washing, boiling, and rinsing the laundry from a pump, well, or spring. Water for the laundry would be hand-carried, heated on a fire for washing, and then poured into a tub. This meant the amount of warm, soapy water was limited; it would be reused to wash the least soiled clothing, then to wash progressively dirtier laundry.
Removal of soap and water from the clothing after washing was a separate process. First, soap would be rinsed out with clear water. After rinsing, the soaking wet clothing would be formed into a roll and twisted by hand to extract water. The entire process often occupied an entire day of work, plus drying and ironing.
Early machines
An early example of washing by machine is the practice of fulling. In a fulling mill, the cloth was beaten with wooden hammers, known as fulling stocks or fulling hammers.
The first English patent under the category of washing machines was issued in 1691. A drawing of an early washing machine appeared in the January 1752 issue of The Gentleman's Magazine, a British publication. Jacob Christian Schäffer's washing machine design was published in 1767 in Germany. In 1782, Henry Sidgier was issued a British patent for a rotating drum washer, and in the 1790s, Edward Beetham sold numerous "patent washing mills" in England.
One of the first innovations in washing machine technology was the use of enclosed containers or basins that had grooves, fingers, or paddles to help with the scrubbing and rubbing of the clothes. The person using the washer would use a stick to press and rotate the clothes along the textured sides of the basin or container, agitating the clothes to remove dirt and mud. This crude agitator technology was hand-powered, but still more effective than actually hand-washing the clothes.
More advancements were made to washing machine technology in the form of the rotating drum design. These early design patents consisted of a drum washer that was hand-cranked to make the wooden drums rotate. While the technology was simple enough, it was a milestone in the history of washing machines, as it introduced the idea of "powered" washing drums. As metal drums started to replace the traditional wooden drums, it allowed for the drum to turn above an open fire or an enclosed fire chamber, raising the water temperature for more effective washes.
It was in the nineteenth century that steam power was first used in washing machine designs.
In 1862, a patented "compound rotary washing machine, with rollers for wringing or mangling" by Richard Lansdale of Pendleton, Manchester, was shown at the 1862 London Exhibition.
The first United States Patent, titled "Clothes Washing", was granted to Nathaniel Briggs of New Hampshire in 1797. Because of the Patent Office fire in 1836, no description of the device survives. The invention of the washing machine is also attributed to Watervliet Shaker Village, as a patent was issued to an Amos Larcom of Watervliet, New York, in 1829, but it is not certain that Larcom was a Shaker. A device that combined a washing machine with a wringer mechanism appeared in 1843 when Canadian John E. Turnbull of Saint John, New Brunswick patented a "Clothes Washer With Wringer Rolls". During the 1850s, Nicholas Bennett of the Mount Lebanon Shaker Society at New Lebanon, New York, invented a "wash mill", but in 1858 he assigned the patent to David Parker of the Canterbury Shaker Village, where it was registered as the "Improved Washing Machine".
Margaret Colvin improved the Triumph Rotary Washer, which was exhibited in the Women's Pavilion at the Centennial International Exhibition of 1876 in Philadelphia. At the same exhibition, the Shakers won a gold medal for their machine.
Electric washing machines were advertised and discussed in newspapers as early as 1904. Alva J. Fisher has been incorrectly credited with the invention of the electric washer. The US Patent Office shows at least one patent issued before Fisher's US patent number 966677 (e.g. Woodrow's US patent number 921195). The first inventor of the electric washing machine remains unknown.
US electric washing machine sales reached 913,000 units in 1928. However, high unemployment rates in the Depression years reduced sales; by 1932 the number of units shipped was down to about 600,000.
An early laundromat in the United States opened in Fort Worth, Texas, in 1934. It was run by Andrew Klein. Patrons used coin-in-the-slot facilities to rent washing machines. The term "laundromat" can be found in newspapers as early as 1884 and they were widespread during the Depression. England established public washrooms for laundry along with bathhouses throughout the nineteenth century.
Washer design improved during the 1930s. The mechanism was now enclosed within a cabinet, and more attention was paid to electrical and mechanical safety. Spin dryers were introduced to replace the dangerous power mangle/wringers of the day.
By 1940, 60% of the 25,000,000 wired homes in the United States had an electric washing machine. Many of these machines featured a power wringer, although built-in spin dryers were not uncommon.
Automatic machines
Bendix Home Appliances, a subsidiary of Avco, introduced the first domestic automatic washing machine in 1937, having applied for a patent in the same year. Avco had licensed the name from Bendix Corporation, an otherwise unrelated company. In appearance and mechanical detail, this first machine was not unlike the front-loading automatic washers produced today.
Although it included many of today's basic features, the machine lacked any drum suspension and therefore had to be anchored to the floor to prevent "walking". Because of the components required, the machine was also expensive. For instance, the Bendix Home Laundry Service Manual (published November 1, 1946) shows that the drum speed change was facilitated by a 2-speed gearbox built to a heavy-duty standard (not unlike a car automatic gearbox, albeit smaller in size). The timer was also probably costly because miniature electric motors were expensive to produce.
Early automatic washing machines were usually connected to a water supply via temporary slip-on connectors to sink taps. Later, permanent connections to hot and cold water became the norm. Most modern front-loading European machines now only have a cold water connection (called "cold fill") and rely completely on internal electric heaters to raise the water temperature.
Many of the early automatic machines had coin-in-the-slot facilities and were installed in the basement laundry rooms of apartment houses.
World War II and after
After the attack on Pearl Harbor, US domestic washer production was suspended for the duration of World War II in favor of manufacturing war material. However, numerous US appliance manufacturers were permitted to undertake the research and development of washers during the war years. Many took the opportunity to develop automatic machines, realizing that these represented the future of the industry.
A large number of US manufacturers introduced competing automatic machines (mainly of the top-loading type) in the late 1940s and early 1950s. General Electric also introduced its first top-loading automatic model in 1947. This machine had many of the features that are incorporated into modern machines. Another early form of automatic washing machine manufactured by The Hoover Company used cartridges to program different wash cycles. This system, called the "Keymatic", used plastic cartridges with key-like slots and ridges around the edges. The cartridge was inserted into a slot on the machine and a mechanical reader operated the machine accordingly.
Several manufacturers produced semi-automatic machines, requiring the user to intervene at one or two points in the wash cycle. A common semi-automatic type (available from Hoover in the UK until at least the 1970s) included two tubs: one with an agitator or impeller for washing, plus another smaller tub for water extraction or centrifugal rinsing. These machines are still available in some countries such as India.
Since their introduction, automatic washing machines have relied on electromechanical timers to sequence the washing and extraction process. Electromechanical timers consist of a series of cams on a common shaft driven by a small electric motor via a reduction gearbox. At the appropriate time in the wash cycle, each cam actuates a switch to engage or disengage a particular part of the machinery (for example, the drain pump motor). One of the first was invented in 1957 by Winston L. Shelton and Gresham N. Jennings, then both General Electric engineers. The device was granted US Patent 2870278.
On the early electromechanical timers, the motor ran at a constant speed throughout the wash cycle, although the user could truncate parts of the program by manually advancing the control dial. However, by the 1950s demand for greater flexibility in the wash cycle led to the introduction of more sophisticated electrical timers to supplement the electromechanical timer. These newer timers enabled greater variation in functions such as the wash time. With this arrangement, the electric timer motor is periodically switched off to permit the clothing to soak and is only re-energized just before a micro-switch being engaged or disengaged for the next stage of the process. Fully electronic timers did not become widespread until decades later.
Despite the high cost of automatic washers, manufacturers had difficulty meeting the demand. Although there were material shortages during the Korean War, by 1953 automatic washing machine sales in the US exceeded those of wringer-type electric machines.
In the UK and most of Europe, electric washing machines did not become popular until the 1950s. This was largely because of the economic impact of World War II on the consumer market, which did not properly recover until the late 1950s. The early electric washers were single-tub wringer-type machines, as fully automatic washing machines were expensive.
During the 1960s, twin tub machines briefly became popular, helped by the low price of the Rolls Razor washers. Twin tub washing machines have two tubs, one larger than the other. The smaller tub in reality is a spinning drum for centrifugal drying while the larger tub only has an agitator in its bottom. Some machines could pump used wash water into a separate tub for temporary storage and to later pump it back for re-use. This was done not to save water or soap, but because heated water was expensive and time-consuming to produce. Automatic washing machines did not become dominant in the UK until well into the 1970s and by then were almost exclusively of the front-loader design.
In early automatic washing machines, any changes in impeller/drum speed were achieved by mechanical means or by a rheostat on the motor power supply. However, since the 1970s electronic control of motor speed has become a common feature on the more expensive models.
Cost-cutting and contemporary development
Over time manufacturers of automatic washers have gone to great lengths to reduce costs. For instance, expensive gearboxes are no longer required, since motor speed can be controlled electronically. Some models can be controlled via WiFi, and have angled/tilted drums to facilitate loading.
Even on some expensive washers, the outer drum of front-loading machines is often (but not always) made of plastic (it can also be made out of metal, but this is expensive). This makes changing the main bearings difficult, as the plastic drum usually cannot be separated into two halves to enable the inner drum to be removed to gain access to the bearing.
Many residential front-loading washing machines typically have a concrete block to dampen vibration. Alternatives include a plastic counterweight that can be filled with water after delivery, reducing or controlling motor speeds, using hydraulic suspensions instead of spring suspensions, and having freely moving steel balls or liquid contained inside a ring mounted on both the top and bottom of the drum to counter the weight of the clothes and reduce vibration.
Most newer front-load machines now use a brushless DC (BLDC) motor directly connected to the basket (direct drive), where the stator assembly is attached to the rear of the outer plastic drum assembly, whilst the co-axial rotor is mounted on the shaft of the inner drum. The direct drive motor eliminates the need for a pulley, belt, and belt tensioner. It was first introduced to washing machines by Fisher and Paykel in 1991. Since then, other manufacturers have followed suit. Some washing machines with this type of motor now come with 10-year or 20-year warranties. The motor type used is an outrunner, due to its slim design with variable speed and high torque. The rotor is connected to the inner tub through its center. It can be made of metal or plastic. Some direct drive washers use induction motors instead of BLDC motors.
Additional features
The modern washing machine market has seen several innovations and features, examples including:
Washing machines including water jets (also known as water sprays, jet sprays and water showers) and steam nozzles that claim to sanitize clothes, help reduce washing times, and remove soil from the clothes. Water jets get their water from the bottom of the drum, thus recirculating the water in the washer.
Others have special drums with holes that will fill with water from the bottom of the tub and redeposit the water on top of the clothes. Some drums have elements with the shape of waves, pyramids, hexagons, domes, or diamonds.
Some include titanium or ceramic heating elements that claim to eliminate calcium buildup in the element. They can heat water up to .
Some high-end models have lights built into the washer itself to light the drum,
Others have soap dispensers where the user fills a tank with detergent and softener and the washing machine automatically doses the detergent and softener and, in some cases, chooses the most appropriate wash cycle. In some models, the tanks come pre-filled and are installed and replaced with new tanks, also pre-filled or refilled by the user, in a dedicated compartment on the bottom of the machine.
Some have support for single-use capsules containing enough laundry additives for one load. The capsules are installed in the detergent compartment.
Many dilute the detergent before it comes in contact with the clothes, some by means of mixing the soap and water with air to make foam, which is then introduced into the drum and improves cleaning performance. Alternatively micro bubbles may be used instead.
Some have pulsators that are mounted on a plate on the bottom of the drum instead of an agitator. The plate spins, and the pulsators generate waves that help shake the soil out of the clothes. Many also include mechanisms to prevent or remove undissolved detergent residue on the detergent dispenser.
It is possible to incorporate a blower and a nozzle to smooth wrinkles in clothes without removing them from the washer.
Some manufacturers like LG Electronics and Samsung Electronics have introduced functions on their washers that allow users to troubleshoot common problems with their washers without having to contact technical support. LG's approach involves a phone receiving signals through sound tones, while Samsung's approach involves having the user take a photo of the washer's time display with a phone. In both methods, the problem and steps to resolve it are displayed on the phone itself. Some models are also NFC enabled. Some implementations are patented under US Patent US20050268669A1 and US Patent US20050097927A1.
In the early 1990s, upmarket machines incorporated microcontrollers for the timing process. These proved reliable and cost-effective, so many cheaper machines now also incorporate microcontrollers rather than electromechanical timers. Since the 2010s, some machines have had touchscreen displays, full-color or color displays, or touch-sensitive control panels.
In 1994, Staber Industries released the System 2000 washing machine, which is the only top-loading, horizontal-axis washer to be manufactured in the United States. The hexagonal tub spins like a front-loading machine, using only about one-third as much water as conventional top-loaders. This factor has led to an Energy Star rating for its high efficiency. This type of horizontal-axis washer and dryer (with a circular drum) is often used in Europe, where space is limited, as they can be as thin as in width.
In 1998, New Zealand-based company Fisher & Paykel introduced its SmartDrive washing machine line in the US. This washing machine uses a computer-controlled system to determine factors such as load size and adjusts the wash cycle to match. It also used a mixed system of washing, first with the "Eco-Active" wash, using a low level of recirculated water being sprayed on the load followed by a more traditional style wash. The SmartDrive also included a direct drive brushless DC electric motor, which simplified the bowl and agitator drive by eliminating the gearbox system.
In 2000, the British inventor James Dyson launched the CR01 ContraRotator, a type of washing machine with two cylinders rotating in opposite directions. It was claimed that this design reduced the wash time and produced cleaner washing than a single-cylinder machine. In 2004 the launch of the CR02, was the first washing machine to gain the British Allergy Foundation Seal of Approval. However, neither of the ContraRotator machines is now in production as they were expensive to manufacture. They were discontinued in 2005. It is patented under U.S. Patent US7750531B2, U.S. Patent US6311527, U.S. Patent US20010023513, U.S. Patent US6311527B1, U.S. Patent USD450164.
In 2001, Whirlpool Corporation introduced the Calypso, the first vertical-axis high-efficiency washing machine to be top-loading. A washplate in the bottom of the tub nutated (a special wobbling motion) to bounce, shake, and toss the laundry. Simultaneously, water containing detergent was sprayed onto the laundry. The machine proved to be good at cleaning but gained a bad reputation due to frequent breakdowns and destruction of laundry. The washer was recalled with a class-action lawsuit and pulled off the market.
In 2003, Maytag introduced their top-loading Neptune TL FAV6800A and TL FAV9800A washers. Instead of an agitator, the machine had two washplates, perpendicular to each other and at a 45-degree angle from the bottom of the tub. The machine would fill with only a small amount of water and the two wash plates would spin, tumbling the load within it, mimicking the action of a front-loading washer in a vertical-axis design.
In 2006, Sanyo introduced the "world-first" (as of February 2, 2006, with regards to home use drum-type washer/dryer) drum-type washing machine with "Air Wash" function (i.e.: using ozone as a disinfectant). It also reused and disinfected rinse water. This washing machine uses only of water in the recycle mode.
Approximately in 2012, eco-indicators were introduced, capable of predicting the energy demand based on the customer settings in terms of program and temperature.
Features available in most modern consumer washing machines:
Delayed execution: a timer to delay the start of the laundry cycle
Predefined programs for different laundry types
Rotation speed settings
Variable temperatures, including cold wash
Additionally, some modern machines feature:
Child lock
Steam
Time remaining indication
Extra water/rinse.
UV disinfection.
Around 2015 and 2017, some manufacturers (namely Samsung and LG Electronics) offered washers and dryers that either have a top-loading washer and dryer built on top of a front-loading washer and dryer respectively (in Samsung washers and dryers) or offer users an optional top-loading washer that can be installed under a washer or dryer (for LG washers and dryers) Both manufacturers have also introduced front-loading washers allowing users to add items after a wash cycle has started, and Samsung has also introduced top-loading washers with a built-in sink and a detergent dispenser that claims to leave no residue on the dispenser itself. In IFA 2017, Samsung released the QuickDrive, a front-loading washer similar to the Dyson ContraRotator but instead of two counter-rotating drums, the QuickDrive has a single drum with a counter-rotating impeller mounted on the back of the drum. Samsung claims this technique reduces cycle times by half and energy consumption by 20%. The US has introduced standards for washing machines that improve their energy efficiency and reduce their water consumption.
Types
Top-loading
The top-loading, vertical-axis washer has been the dominant design in the United States and Canada. This design places the clothes in a vertically mounted perforated basket that is contained within a water-retaining tub, with a finned water-pumping agitator in the center of the bottom of the basket. Clothes are loaded through the top of the machine, which is usually but not always covered with a hinged door. The drum of a top loading washing machine can include a lint trap.
Agitation
During the wash cycle, the outer tub is filled with water sufficient to fully immerse and suspend the clothing freely in the basket. The movement of the agitator pushes water outward between the paddles towards the edge of the tub. The water then moves outward, up the sides of the basket, towards the center, and then down towards the agitator to repeat the process, in a circulation pattern similar to the shape of a torus. The agitator direction is periodically reversed because continuous motion in one direction would just lead to the water spinning around the basket with the agitator rather than the water being pumped in the torus-shaped motion. Some washers supplement the water-pumping action of the agitator with a large rotating screw on the shaft above the agitator, to help move water downwards in the center of the basket. A washing machine can have an impeller, also called a wash plate, instead of an agitator, which serves the same purpose but does not have a vertical cylinder extending from its base.
Since the agitator and the drum are separate and distinct in a top-loading washing machine, the mechanism of a top-loader is inherently more complicated than a front-loading machine. Manufacturers have devised several ways to control the motion of the agitator during the wash and rinse separately from the high-speed rotation of the drum required for the spin cycle. While a top-loading washing machine could use a universal motor or DC brushless motor, it is conventional for top-loading washing machines to use more expensive, heavy, and potentially more electrically efficient and reliable induction motors.
An alternative to this oscillating agitator design is the impeller-type washtub pioneered by Hoover on its long-running Hoovermatic series of top-loading machines. Here, an impeller (trademarked by Hoover as a "Pulsator") mounted on the side of the tub spins in a constant direction and creates a fast-moving current of water in the tub which drags the clothes through the water along a toroidal path. This design was used in the Hoover 0307 washer. The impeller design has the advantage of mechanical simplicity – a single-speed motor with belt drive is all that is required to drive the Pulsator with no need for gearboxes or complex electrical controls, but has the disadvantage of lower load capacity in relation to tub size. Hoovermatic machines were made mostly in twin-tub format for the European market (where they competed with Hotpoint's Supermatic line which used the oscillating agitator design) until the early 1990s. Some industrial garment testing machines still use the Hoover wash action. Another alternative involves 'pulsating' the agitator, in other words having an agitator with a reciprocating motion along its vertical axis. Some washing machines have agitators that move in an orbiting motion or agitators that nutate at the bottom. Special top loading washing machines designed for washing sneakers can incorporate bristles in their agitators. Alternatively the inner tub itself can nutate inside the outer tub.
The many different ways manufacturers have solved the same problem over the years is a good example of many different ways to solve the same engineering problem with different goals, different manufacturing capabilities and expertise, and different patent encumbrances.
Reversible motor
In many current top-loading washers, if the motor spins in one direction, the gearbox drives the agitator; if the motor spins the other way, the gearbox locks the agitator and spins the basket and agitator together. Similarly, if the pump motor rotates one way it recirculates the sudsy water; in the other direction it pumps water from the machine during the spin cycle. Mechanically, this system is very simple.
Mode-changing transmission
In some top-loaders, the motor runs only in one direction. During agitation, the transmission converts the rotation into the alternating motion driving the agitator. During the spin cycle, the timer turns on a solenoid which engages a clutch locking the motor's rotation to the wash basket, providing a spin cycle. General Electric's very popular line of Filter-Flo (seen to the right) used a variant of this design where the motor reversed only to pump water out of the machine. The same clutch which allows the heavy tub full of wet clothes to "slip" as it comes up to the motor's speed, is also allowed to "slip" during agitation to engage a Gentle Cycle for delicate clothes.
Whirlpool (Kenmore) created a popular design demonstrating the complex mechanisms which could be used to produce different motions from a single motor with the so-called "wig wag" mechanism, which was used for decades until modern controls rendered it obsolete. In the Whirlpool mechanism, a protruding moving piece oscillates in time with the agitation motion. Two solenoids are mounted to this protruding moving piece, with wires attaching them to the timer. During the cycle, the motor operates continuously, and the solenoids on the "wig wag" engage in agitation or spin. Despite the wires controlling the solenoids being subject to abrasion and broken connections due to their constant motion and the solenoids operating in a damp environment where corrosion could damage them, these machines were surprisingly reliable.
Reversible motor with mode-changing transmission
Some top-loaders, especially compact apartment-sized washers, use a hybrid mechanism. The motor reverses direction every few seconds, often with a pause between direction changes, to perform the agitation. The spin cycle is accomplished by engaging a clutch in the transmission. A separate motorized pump is generally used to drain this style of machine. These machines could easily be implemented with universal motors or more modern DC brushless motors, but older ones tend to use a capacitor-start induction motor with a pause between reversals of agitation.
Front-loading
The front-loading or horizontal-axis clothes washer is the dominant design in Europe and in most parts of the world. In the United States and Canada, most "high-end" washing machines are of this type. In addition, most commercial and industrial clothes washers around the world are of the horizontal-axis design.
This layout mounts the inner drum and outer drum horizontally, and loading is through a door at the front of the machine. The door often but not always contains a transparent window. Agitation is supplied by the back-and-forth rotation of the cylinder and by gravity. The clothes are lifted by paddles on the inside wall of the drum and then dropped. This motion flexes the weave of the fabric and forces water and detergent solution through the clothes load. Because the wash action does not require the clothing to be freely suspended in water, only enough water is needed to moisten the fabric. Because less water is required, front-loaders typically use less soap, and the repeated dropping and folding action of the tumbling can easily produce large amounts of foam or suds.
Front-loaders control water usage through the surface tension of water, and the capillary wicking action this creates in the fabric weave. A front-loader washer always fills to the same low water level, but a large pile of dry clothing standing in water will soak up the moisture, causing the water level to drop. The washer then refills to maintain the original water level. Because it takes time for this water absorption to occur with a motionless pile of fabric, nearly all front-loaders begin the washing process by slowly tumbling the clothing under the stream of water entering and filling the drum, to rapidly saturate the clothes with water.
Compared to top-loading washers, clothing can be packed more tightly in a front loader, up to the full drum volume if using a cotton wash cycle. This is because wet cloth usually fits into a smaller space than dry cloth, and front-loaders can self-regulate the water needed to achieve correct washing and rinsing. However, extreme overloading of front-loading washers pushes fabrics towards the small gap between the loading door and the front of the wash basket, potentially resulting in fabrics lost between the basket and outer tub, and in severe cases, tearing of clothing and jamming the motion of the basket.
Mechanical aspects
Front-loading washers are mechanically simple compared to top-loaders, with the main motor (a universal motor or variable-frequency drive motor) normally being connected to the drum via a grooved pulley belt and large pulley wheel without the need for a gearbox, clutch or crank. The action of a front-loading washing machine is better suited to a motor capable of reversing direction with every reversal of the wash drum; a universal motor is noisier, less efficient, and does not last as long, but is better suited to the task of reversing direction every few seconds. Some models, such as those by LG, use a motor directly connected to the drum, eliminating the need for a belt and pulley.
However, front-load washers suffer from their own technical challenges due to the horizontal disposition of the drum. A top-loading washer keeps water inside the tub merely through the force of gravity pulling down on the water, while a front-loader must tightly seal the door with a gasket to prevent water dripping onto the floor during the wash cycle. This access door is locked shut with an interlocking device during the entire wash cycle, since opening the door with the machine in use could result in water gushing onto the floor. If this interlock is broken for any reason, such a machine stops operation, even if this failure happens mid-cycle. In most machines, the interlock is usually doubly redundant to prevent either opening with the drum full of water or being opened during the spin cycle. For front-loaders without viewing windows on the door, it is possible to accidentally pinch the fabric between the door and the drum, resulting in tearing and damage to the pinched clothing during tumbling and spinning.
Nearly all front-loader washers for the consumer market also use a folded flexible bellows assembly around the door opening to keep clothing contained inside the drum during the tumbling wash cycle. If this bellows assembly were not used, small articles of clothing such as socks could slip out of the wash drum near the door and fall down the narrow slot between the outer and inner drums, plugging the drain and possibly jamming rotation of the inner drum. Retrieving lost items from between the outer drum and inner drum can require complete disassembly of the front of the washer and pulling out the entire inner wash drum. Commercial and industrial front-loaders used by businesses (described below) usually do not use the bellows, but instead require all small objects to be placed in a mesh bag to prevent loss near the drum opening.
Variant and hybrid designs
There are many variations of the two general designs. Top-loading machines in Asia use impellers instead of agitators. Impellers are similar to agitators except that they do not have the center post extending up in the middle of the washtub basket.
Horizontal-axis top-loader
Some machines which load from the top are otherwise much more similar to front-loading horizontal-axis drum machines. They have a drum rotating around a horizontal axis, as a front-loader, but there is no front door; instead, there is a liftable lid that provides access to the drum, which has a hatch that can be latched shut. Clothes are loaded, the hatch and lid are closed, and the machine operates and spins just like a front loader. These machines are narrower but usually taller than front-loaders, usually have a lower capacity, and are intended for use where only a narrow space is available, as is sometimes the case in Europe. They have incidental advantages: they can be loaded while standing (but force the user to bend down instead of crouching down or sitting to unload); they do not require a perishable rubber bellows seal; and instead of the drum having a single bearing on one side, it has a pair of symmetrical bearings, one on each side, avoiding asymmetrical bearing loading and potentially increasing life.
Combo washer dryer
There are also combo washer dryer machines that combine washing cycles and a full drying cycle in the same drum, eliminating the need to transfer wet clothes from a washer to a dryer machine. In principle, these machines are convenient for overnight cleaning (the combined cycle is considerably longer), but the effective capacity for cleaning larger batches of laundry is drastically reduced. The drying process tends to use much more energy than using two separate devices, because a combo washer dryer not only must dry the clothing but also needs to dry out the wash chamber itself.
These machines are used more where space is at a premium, such as areas of Europe and Japan because they can be fit into small spaces, perform both washing and drying, and many can be operated without dedicated utility connections. In these machines, the washer and dryer functions often have different capacities, with the dryer usually having the lowest capacity.
These combo machines should not be confused with a dryer on top of a washer installation, or with a laundry center, which is a one-piece appliance offering a compromise between a washer-dryer combo and a full washer to the side of the dryer installation or a dryer on top of a washer installation. Laundry centers usually have the dryer on top of the washer, with the controls for both machines being on a single control panel. Often, the controls are simpler than the controls on a washer-dryer combo or a dedicated washer and dryer. Some implementations are patented under US Patent US6343492B1 and US Patent US 6363756B1.
Comparison
True front-loading machines, top-loading machines with horizontal-axis drums, and true top-loading vertical-axis machines can be compared on several aspects:
Efficient cleaning: Front loaders usually use less energy, water, and detergent compared to the best top-loaders. High-efficiency washers use 20% to 60% of the detergent, water, and energy of "standard" commonly-used top-loader washers. They usually take somewhat longer (20–110 minutes) to wash a load, but are often computer controlled with additional sensors, to adapt the wash cycle to the needs of each load.
Water usage: Front-loaders usually use less water than top-loading residential clothes washers. Estimates are that front-loaders use from one-third to one half as much water as top-loaders.
Spin-dry effectiveness: Front-loaders (and European horizontal-axis top-loaders and some front-loaders) offer much higher maximum spin speeds of up to 2000 RPM, although home machines tend to be in the 1000 to 1400 RPM range, while top-loaders (with agitators) do not exceed 1140 RPM. High-efficiency top-loaders with a wash plate (instead of an agitator) can spin up to 1100 RPM, as their center of gravity is lower. Higher spin speeds, along with the diameter of the drum, determine the g-force, and a higher g-force removes more residual water, making clothes dry faster. This also reduces energy consumption if clothes are dried in a clothes dryer.
Cycle length: Top-loaders have tended to have shorter cycle times, in part because their design has traditionally emphasized simplicity and speed of operation more than resource conservation. It is observed that top-loaders wash the clothes in half the time as compared to a front-load washing machine.
Wear and abrasion: Top-loaders require an agitator or impeller mechanism to force enough water through clothes to clean them effectively, which greatly increases mechanical wear and tear on fabrics. Front-loaders use paddles in the drum to repeatedly pick up and drop clothes into the water for cleaning; this gentler action causes less wear and tear. The rate of clothes wear can be roughly gauged by the amount of accumulation in a clothes dryer lint filter, since the lint largely consists of stray fibers detached from textiles during washing and drying.
Difficult items: Top-loaders may have trouble cleaning large items, such as sleeping bags or pillows, which tend to float on top of the wash water rather than circulate within it. In addition, vigorous top-loader agitator motions may damage delicate fabrics. Whereas in a front-load washing machine, one can easily wash pillows, shoes, soft toys, and other difficult-to-wash items.
Noise: Front-loaders tend to operate more quietly than top-loaders because the door seal helps contain noise, and because there is less of a tendency towards imbalance. Top loaders usually need a mechanical transmission (due to agitators, see above), which can generate more noise than the rubber belt or direct drive found in most front-loaders.
Compactness: True front-loading machines may be installed underneath counter-height work surfaces. A front-loading washing machine, in a fully fitted kitchen, may even be disguised as a kitchen cabinet. These models can also be convenient in homes with limited floor area, since the clothes dryer may be installed directly above the washer ("stacked" configuration).
Water leakage: Top-loading machines are less prone to leakage because simple gravity reliably keeps water from spilling out the loading door on top. True front-loading machines require a flexible seal or gasket on the front door, and the front door must be locked during operation to prevent opening, lest large amounts of water spill out. This seal may leak and require replacement. However, many current front-loaders use so little water that they can be stopped mid-cycle for the addition or removal of laundry, while keeping the water level in the horizontal tub below the door level. Best practice installations of either type of machine will include a floor drain or an overflow catch tray with a drain connection, since neither design is immune to leakage or a solenoid valve getting stuck in the open position.
Maintenance and reliability: Top-loading washers are more tolerant of maintenance neglect, and may not need a regular "freshening" cycle to clean door seals and bellows. During the spin cycle, a top-loading tub is free to move about inside the cabinet of the machine, using only a lip around the top of the inner basket and outer tub to keep the spinning water and clothing from spraying out over the edge. Therefore, the potentially problematic door-sealing and door-locking mechanisms used by true front-loaders are not needed. On the other hand, top-loaders use mechanical gearboxes that are more vulnerable to wear than simpler front-load motor drives.
Accessibility and ergonomics: Front-loaders are more convenient for shorter people and those with paraplegia, as the controls are front-mounted and the horizontal drum eliminates the need for standing or climbing. Risers, also referred to as pedestals, often with storage drawers underneath, can be used to raise the door of a true front-loader closer to the user's level. However, if stacked, the dryer controls, if at the top of the dryer, may be too tall for shorter people to conveniently access.
Initial cost: In countries where top-loaders are popular, front-loaders tend to be more expensive to buy than top-loaders, though their lower operating costs can lead to lower total cost of ownership, especially if energy, detergent, or water are expensive. On the other hand, in countries with a large front-loader user base, top-loaders are usually seen as alternatives and more expensive than basic off-brand front-loaders, although without many differences in total cost of ownership apart from design-originated ones. In addition, manufacturers have tended to include more advanced features such as internal water heating, automatic dirt sensors, and high-speed emptying on front loaders, although some of these features could be implemented on top loaders.
Wash cycles
The earliest washing machines simply carried out a washing action when loaded with clothes and soap, filled with hot water, and started. Over time machines became more and more automated, first with complex electromechanical controllers, then fully electronic controllers; users put clothes into the machine, select a suitable program via a switch, start the machine, and come back to remove clean and slightly damp clothes at the end of the cycle. The controller starts and stops many different processes including pumps and valves to fill and empty the drum with water, heating, and rotating at different speeds, with different combinations of settings for different fabrics.
Longer wash cycles can allow greater water and energy efficiency (with less water to heat up). For a load, from 2011 to 2021, the average Australian washing machine cycle (including rinsing and spinning) has lengthened from 99 to 144 minutes for front-loaders, and 55 to 59 minutes for top-loaders.
Washing
Many front-loading machines have internal electrical heating elements to heat the wash water, to near boiling if desired. The rate of the chemical cleaning action of the detergent and other laundry chemicals increases greatly with temperature, by the Arrhenius equation. Washing machines with internal heaters can use special detergents formulated to release different chemical ingredients at different temperatures, allowing different types of stains and soils to be cleaned from the clothes as the wash water is heated by the electrical heater.
However, higher-temperature washing uses more energy, and many fabrics and elastics are damaged at higher temperatures. Temperatures exceeding have the undesirable effect of deactivating the enzymes when using biological detergent.
Many machines are cold-fill, connected to cold water only, which they internally heat to operating temperature. Where water can be heated more cheaply or with less carbon dioxide emission than by electricity, a cold-fill operation is inefficient.
Front-loaders need to use low-sudsing detergents because the tumbling action of the drum entrains air into the clothes load, which can cause excessive foamy suds and overflows. However, due to the efficient use of water and detergent, the suds issue with front-loaders can be controlled by simply using less detergent, without lessening the cleaning action.
Rinsing
Washing machines perform several rinses after the main wash to remove most of the detergent. Modern washing machines use less hot water due to environmental concerns; however, this has led to the problem of poor rinsing on many washing machines on the market, which can be a problem to people who are sensitive to detergents. The Allergy UK website suggests re-running the rinse cycle, or rerunning the entire wash cycle without detergent.
In response to complaints, many washing machines allow the user to select additional rinse cycles, at the expense of higher water usage and longer cycle time. Bosch, for example, in its allergy wash program, incorporates an additional three-minute rinse cycle with water of at least to rinse off detergent residues and any allergens.
Spin
Front-loading machines spin in multiple stages of their cycle: after main wash, after individual rinses, and the final high-speed spin. Some of those spins may be absent depending on the particular cycle.
Higher spin speeds, along with larger tub diameters, remove more water, leading to faster drying. On the other hand, the need for ironing can be reduced by not using the spin cycle in the washing machine.
If a heated clothes dryer is used after the wash and spin, energy use is reduced if more water has been removed from clothes. However, faster spinning can crease clothes more. Also, mechanical wear on bearings increases rapidly with rotational speed, reducing life. Early machines would spin at 300 rpm and, because of lack of any mechanical suspension, would often shake and vibrate.
In 1976, most front-loading washing machines spun at around 700 RPM, or less. Today, most machines spin at 1000–1600 RPM. Most machines have variable speeds, ranging 300–2000 RPM depending on the machine.
Separate spin-driers, without washing functionality, are available for specialized applications. For example, a small high-speed centrifuge machine may be provided in locker rooms of communal swimming pools to allow wet swimsuits to be substantially dried to a slightly damp condition after daily use.
Washing machines often incorporate balance rings filled with a liquid such as a calcium chloride salt water solution, that are designed to balance the inner drum of the washer during spin cycles. The balance ring may be filled with oil and contain balls on races, somewhat similarly to a ball bearing, to achieve the same effect. The Bendix Economat used a flexible rubber inner tub that would squeeze the clothes towards the agitator located in the center of the inner tub in order to remove water from the clothes, instead of spinning the inner tub. This was performed by exerting a vacuum on the inner tub.
Maintenance wash
Many home washing machines use a plastic, rather than metal, outer shell to contain the wash water; residue can build up on the plastic tub over time. Some manufacturers advise users to perform a regular maintenance or "freshening" wash to clean the inside of the washing machine of any mold, bacteria, encrusted detergent, and unspecified dirt more effectively than with a normal wash.
A maintenance wash is performed without any laundry, on the hottest wash program, adding substances such as white vinegar, 100 grams of citric acid, a detergent with bleaching properties, or a proprietary washing machine cleaner. The first injection of water goes into the sump so the machine can be allowed to fill for about 30 seconds before adding cleaning substances.
Installation and flood prevention
Flexible rubber hoses are typically used to connect from a building water supply to a washing machine. These hoses are often exposed to full water pressure on a continuing basis and can deteriorate over time, developing bulges or weak spots that eventually cause leaks or catastrophic bursting and flooding. Since the hoses are often hidden from view, they may be difficult to inspect and easily forgotten until a problem occurs. If a hose burst occurs when nobody is present to notice the problem, a huge volume of water can be delivered over a short time, causing extensive interior flooding damage or even structural damage. It has been estimated that a burst supply hose can deliver two tons of water in an hour.
To reduce these risks, it is a common recommendation to use flexible hoses which have been jacketed with a braided stainless steel mesh. This jacketing cannot prevent leaks from developing, but it can slow the development of large bulges or "aneurysms" which can burst suddenly without warning. However, even braided metal jackets often cannot withstand the enormous pressures generated by water freezing within an enclosed volume.
An additional precaution is to install a washing machine inside a shallow metal or plastic pan, which can collect minor leakage and divert the water to a nearby drain, or to the outside of a building. Drain pans can also divert water released by other problems, such as a jammed solenoid valve in a washing machine. A serious limitation of drain pans is that they typically cannot handle the large volumes of pressurized water released by a burst supply hose, so a drain pan is no substitute for hose burst precautions. In the absence of a drain, a pan may still be useful to confine leakage temporarily, while a local or remote water alarm is triggered.
In addition or instead of an alarm, a water detector may signal the main water shutoff valve to the building to be automatically closed to prevent flooding.
A very effective precaution is to install a shutoff or isolation valve which stops any water from being supplied, except when a washing machine is actually operating. The simplest method is to manually open and close the hot and cold water shutoff valves (traditionally globe valves) behind the washing machine, each time it is used. This method relies on the washing machine user conscientiously operating the two valves each time laundry is done, in spite of the awkward location of the valves and the tedious process of turning the handles through multiple rotations.
An improvement over the traditional setup is to install a specialized laundry shutoff valve. Typically, it consists of two ball valves connected to a single handle, so they can be operated by a horizontal or vertical lever moved by 90 degrees. This makes the operation of the valves a quick procedure, but the washing machine user must still remember to turn off the water, even though the failure to do this produces no immediately obvious problems.
To close this risk exposure, some shutoff valves have a spring-energized mechanical timer which is started when the user pushes a lever to open the valves. After a preset time of several hours elapses, the spring-powered mechanism automatically closes the valve without further user intervention. A variant of this setup requires the user to press a button to open the valves for an electrically-timed interval.
Other automatic valve operating mechanisms electronically detect when a washing machine draws electrical power as it starts, and then open the water supply valves. Typically, the power plug for the washing machine is connected to a special detector receptacle or cable, to allowing monitoring of the power draw.
Although pressurized water supply leaks can cause the most damage in the least amount of time, water drainage can also cause problems if not handled properly. Washing machine drainage hoses should be secured properly to prevent accidental dislodgement, and drains should be inspected and cleared periodically to prevent buildup of laundry lint, mold, and other deposits.
Efficiency and standards
Capacity and cost are both considerations when purchasing a washing machine. All else being equal, a machine of higher capacity will cost more to buy, but will be more convenient if large amounts of laundry must be cleaned. Fewer runs of a machine of larger capacity may have lower running costs and better energy and water efficiency than frequent use of a smaller machine, particularly for large families. However, running a large machine with small loads is typically inefficient and wasteful, unless the machine has been designed to handle such situations.
For many years energy and water efficiency were not regulated, and little attention was paid to them. From the last part of the 20th century, increasing attention was paid to efficiency, with regulations enforcing some standards. Efficiency became a selling point, both to save on running costs and to reduce carbon dioxide emissions associated with energy generation, and waste of water.
As energy and water efficiency became regulated, they became a selling point for buyers; however, the effectiveness of rinsing was not specified, and it did not directly attract the attention of buyers. Therefore, manufacturers tended to reduce the degree of rinsing after washing, saving water and electrical energy. This had the side-effect of leaving more detergent residue in clothes, which can affect people with allergies or sensitivity. In response to complaints, some manufacturers have now designed their machines with a user-selectable option for additional rinsing.
Europe
Washing machines display an EU Energy Label with grades for energy efficiency, washing performance, and spin efficiency. Grades for energy efficiency run from A+++ to D (best to worst), providing a simple method for judging running costs. Washing performance and spin efficiency are graded in the range A to G. However, all machines for sale must have washing performance A, so that manufacturers cannot compromise washing performance in order to improve the energy efficiency. This labeling has had the desired effect of driving customers toward more efficient washing machines and away from less efficient ones.
According to regulations, each washing machine is equipped with a wastewater filter. This ensures that no hazardous chemical substances are disposed of improperly through the sewage system; on the other hand, it also ensures that if there is backflow in the plumbing system, sewage cannot enter the washing machine.
United States
Top-loading and front-loading clothes washers are covered by a single national standard regulating energy consumption. The old federal standards applicable before January 2011 did not restrict water consumption; there was no limit on how much unheated rinse water could be used. Energy consumption for clothes washers is quantified using the energy factor.
After new mandatory federal standards were introduced, many US washers were manufactured to be more energy- and water-efficient than required by the federal standard, or even than required by the more-stringent Energy Star standard. Manufacturers were further motivated to exceed mandatory standards by a program of direct-to-manufacturer tax credits.
In North America, the Energy Star program compares and lists energy-efficient clothes washers. Certified Energy Star units can be compared by their Modified Energy Factor (MEF) and Water Factor (WF) coefficients.
The MEF figure of merit states how many cubic feet (about 28.3 liters) of clothes are washed per kWh (kilowatt hour). The coefficient is influenced by factors including the configuration of the washer (top-loading, front-loading), its spin speed, and the temperatures and the amount of water used in the rinse and wash cycles.
Energy Star residential clothes washers must have an MEF of at least 2.0 (the higher the better); the best machines may reach 3.5. Energy Star washers must also have a WF of less than 6.0 (the lower the better).
Commercial use
A commercial washing machine is intended for more intensive use than a consumer washing machine. Durability and functionality is more important than style; most commercial washers are bulky and heavy, often with more expensive stainless steel construction to minimize corrosion in a constantly-moist environment. They are built with large easy-to-open service covers, and washers are designed not to require access from the underside for service. Commercial washers are often installed in long rows, with a wide access passageway behind all the machines to allow maintenance without moving the heavy machinery.
Laundromat machines
Many commercial washers are built for use by the general public, and are installed in publicly accessible laundromats or laundrettes. Originally, they were operated by coins (similar to older vending machines), but today they are activated by money accepting devices or card readers. The features of a commercial laundromat washer are usually more limited than those of a consumer washer, usually offering just two or three basic wash programs and an option to choose wash cycle temperatures. Some more-advanced models allow extra-cost options such as an additional wash or rinse cycle, at the choice of the user.
The typical front-loading commercial washing machine also differs from consumer models in its discharge of spent wash and rinse water. While the consumer models pump used washer water out, allowing the waste drainage pipe to be located above the floor level, front-loading commercial machines generally use only gravity to expel used water. A drain valve at the bottom rear of the machine opens at the appointed time during the cycle, allowing water to flow out. This requires a special drainage trough equipped with a filter and drain, and routed behind each machine. The trough is usually part of a cement platform built for the purpose of raising the machines to a convenient height, and can be seen behind washers at most laundromats.
Most laundromat machines are horizontal-axis front-loading models, because of their lower operating costs (notably, lower consumption of expensive hot water).
Industrial washers
By contrast, commercial washers for internal business operations (which are often referred to as "washer/extractor" machines) may include features absent from domestic machines. Many commercial washers offer an option for automatic injection of five or more different chemical types, so that the operator does not have to deal with constantly measuring out soap products and fabric softeners for each load by hand. Instead, a precise metering system draws the detergents and wash additives directly from large liquid-chemical storage barrels, and injects them as needed into the various wash and rinse cycles. Some computer-controlled commercial washers offer the operator detailed control over the various wash and rinse cycles, allowing the operator to program custom washing cycles.
Most large-scale industrial washers are horizontal-axis machines, but they may have front-, side-, or top-load doors. Some industrial clothes washers can batch-process up to of textiles at once, and can be used for extremely machine-abusive washing tasks such as stone washing or fabric bleaching and dyeing.
An industrial washer can be mounted on heavy-duty shock absorbers and attached to a concrete floor, so that it can extract water from even the most severely out-of-balance and heavy wash loads. Noise and vibration are not as unacceptable as in a domestic machine. The machine may be mounted on hydraulic cylinders, permitting the entire washer to be lifted and tilted so that fabrics can be automatically dumped from the wash drum onto a conveyor belt once the cycle is complete.
One special type of continuous-processing washer is known as the tunnel washer. This specialized high-capacity machine does not have a drum where everything being washed undergoes distinct wash and rinse cycles. Instead, the laundry progresses slowly and continuously through a long, large-diameter horizontal-axis rotating tube in the manner of an assembly line, with different processes at different positions.
Social impact
The historically laborious process of washing clothes (a task which often consumed a whole day) was at times described as "women's work". The spread of the washing machine has been seen to be a force behind the improvement of women's position in society.
Before the advent of the washing machine, laundry was done first at watercourses, and later in public wash-houses known as lavoirs. Camille Paglia and others argue that the washing machine led to a type of social isolation of women, as a previously communal activity became a solitary one.
In 2009 the Italian newspaper L'Osservatore Romano reprinted a Playboy magazine article on International Women's Day arguing that the washing machine had done more for the liberation of women than the contraceptive pill and abortion rights. A study from Université de Montréal, Canada presented a similar point of view, and added refrigerators. The following year, Swedish statistician Hans Rosling suggested that the positive effect the washing machine had on the liberation of women makes it "the greatest invention of the industrial revolution". It has been argued that washing machines are an example of labor-saving technology which does not decrease employment, because households can internalize the gains of the innovation.
Historian Frances Finnegan credits the rise of domestic laundry technology in helping to undercut the economic viability of the Magdalene asylums in Ireland (later revealed to be inhumanly abusive prisons for women), by supplanting their laundry businesses and prompting the eventual closure of the institutions as a whole. Irish feminist Mary Frances McDonald has described washing machines as the single most life-changing invention for women.
In India, dhobis, a caste group specialized in washing clothes, are slowly adapting to modern technology, but even with access to washing machines, many still handwash garments as well. Since most modern homes are equipped with a washing machine, many Indians have dispensed with the services of the dhobiwallahs.
Environmental impact
Due to the increasing cost of repairs relative to the price of a washing machine, there has been a major increase in the yearly number of defective washing machines being discarded, to the detriment of the environment. The cost of repair and the expected life of a machine may make the purchase of a new machine seem like the better option.
Different washing machine models vary widely in their use of water, detergent, and energy. The energy required for heating is large compared to that used by lighting, electric motors, and electronic devices. Because of their use of hot water, washing machines are among the largest overall consumers of energy in a typical modern home.
Washing machines worldwide release around 62 million tonnes of carbon dioxide equivalent in a year. However, modern improvements have been made aiming to lower these emission numbers, and it depends on the user's choice to fully determine their environmental impact.
| Technology | Household appliances | null |
172121 | https://en.wikipedia.org/wiki/Phonograph%20record | Phonograph record | A phonograph record (also known as a gramophone record, especially in British English) or a vinyl record (for later varieties only) is an analog sound storage medium in the form of a flat disc with an inscribed, modulated spiral groove. The groove usually starts near the outside edge and ends near the center of the disc. The stored sound information is made audible by playing the record on a phonograph (or "gramophone", "turntable", or "record player").
Records have been produced in different formats with playing times ranging from a few minutes to around 30 minutes per side. For about half a century, the discs were commonly made from shellac and these records typically ran at a rotational speed of 78 rpm, giving it the nickname "78s" ("seventy-eights"). After the 1940s, "vinyl" records made from polyvinyl chloride (PVC) became standard replacing the old 78s and remain so to this day; they have since been produced in various sizes and speeds, most commonly 7-inch discs played at 45 rpm (typically for singles, also called 45s ("forty-fives")), and 12-inch discs played at 33⅓ rpm (known as an LP, "long-playing records", typically for full-length albums) – the latter being the most prevalent format today.
Overview
The phonograph record was the primary medium used for music reproduction throughout the 20th century. It had co-existed with the phonograph cylinder from the late 1880s and had effectively superseded it by around 1912. Records retained the largest market share even when new formats such as the compact cassette were mass-marketed. By the 1980s, digital media, in the form of the compact disc, had gained a larger market share, and the record left the mainstream in 1991. Since the 1990s, records continue to be manufactured and sold on a smaller scale, and during the 1990s and early 2000s were commonly used by disc jockeys (DJs), especially in dance music genres. They were also listened to by a growing number of audiophiles. The phonograph record has made a niche resurgence in the early 21st century, growing increasingly popular throughout the 2010s and 2020s.
Phonograph records are generally described by their diameter in inches (12-inch, 10-inch, 7-inch), the rotational speed in revolutions per minute (rpm) at which they are played (, , , 45, 78), and their time capacity, determined by their diameter and speed (LP [long play], 12-inch disc, rpm; EP [extended play], 12-inch disc or 7-inch disc, or 45 rpm; Single, 7-inch or 10-inch disc, 45 or 78 rpm); their reproductive quality, or level of fidelity (high-fidelity, orthophonic, full-range, etc.); and the number of audio channels (mono, stereo, quad, etc.).
The phrase broken record refers to a malfunction when the needle skips/jumps back to the previous groove and plays the same section over and over again indefinitely.
Naming
The various names have included phonograph record (American English), gramophone record (British English), record, vinyl, LP (originally a trademark of Columbia Records), black disc, album, and more informally platter, wax, or liquorice pizza.
Early development
Manufacture of disc records began in the late 19th century, at first competing with earlier cylinder records. Price, ease of use and storage made the disc record dominant by the 1910s. The standard format of disc records became known to later generations as "78s" after their playback speed in revolutions per minute, although that speed only became standardized in the late 1920s. In the late 1940s new formats pressed in vinyl, the 45 rpm single and 33 rpm long playing "LP", were introduced, gradually overtaking the formerly standard "78s" over the next decade. The late 1950s saw the introduction of stereophonic sound on commercial discs.
Predecessors
The phonautograph was invented by 1857 by Frenchman Édouard-Léon Scott de Martinville. It could not, however, play back recorded sound, as Scott intended for people to read back the tracings, which he called phonautograms. Prior to this, tuning forks had been used in this way to create direct tracings of the vibrations of sound-producing objects, as by English physicist Thomas Young in 1807.
In 1877, Thomas Edison invented the first phonograph, which etched sound recordings onto phonograph cylinders. Unlike the phonautograph, Edison's phonograph could both record and reproduce sound, via two separate needles, one for each function.
The first disc records
The first commercially sold disc records were created by Emile Berliner in the 1880s. Emile Berliner improved the quality of recordings while his manufacturing associate Eldridge R. Johnson, who owned a machine shop in Camden, New Jersey, eventually improved the mechanism of the gramophone with a spring motor and a speed regulating governor, resulting in a sound quality equal to Edison's cylinders. Abandoning Berliner's "Gramophone" trademark for legal reasons in the United States, Johnson's and Berliner's separate companies reorganized in 1901 to form the Victor Talking Machine Company in Camden, New Jersey, whose products would come to dominate the market for several decades.
Berliner's Montreal factory, which became the Canadian branch of RCA Victor, still exists. There is a dedicated museum in Montreal for Berliner (Musée des ondes Emile Berliner).
78 rpm disc developments
Early speeds
Early disc recordings were produced in a variety of speeds ranging from 60 to 130 rpm, and a variety of sizes. As early as 1894, Emile Berliner's United States Gramophone Company was selling single-sided 7-inch discs with an advertised standard speed of "about 70 rpm".
One standard audio recording handbook describes speed regulators, or governors, as being part of a wave of improvement introduced rapidly after 1897. A picture of a hand-cranked 1898 Berliner Gramophone shows a governor and says that spring drives had replaced hand drives. It notes that:
The speed regulator was furnished with an indicator that showed the speed when the machine was running so that the records, on reproduction, could be revolved at exactly the same speed...The literature does not disclose why 78 rpm was chosen for the phonograph industry, apparently this just happened to be the speed created by one of the early machines and, for no other reason continued to be used.
In 1912, the Gramophone Company set 78 rpm as their recording standard, based on the average of recordings they had been releasing at the time, and started selling players whose governors had a nominal speed of 78 rpm. By 1925, 78 rpm was becoming standardized across the industry. However, the exact speed differed between places with alternating current electricity supply at 60 hertz (cycles per second, Hz) and those at 50 Hz. Where the mains supply was 60 Hz, the actual speed was 78.26 rpm: that of a 60 Hz stroboscope illuminating 92-bar calibration markings. Where it was 50 Hz, it was 77.92 rpm: that of a 50 Hz stroboscope illuminating 77-bar calibration markings.
At least one attempt to lengthen playing time was made in the early 1920s. World Records produced records that played at a constant linear velocity, controlled by Noel Pemberton Billing's patented add-on speed governor.
Acoustic recording
Early recordings were made entirely acoustically, the sound was collected by a horn and piped to a diaphragm, which vibrated the cutting stylus. Sensitivity and frequency range were poor, and frequency response was irregular, giving acoustic recordings an instantly recognizable tonal quality. A singer almost had to put their face in the recording horn. A way of reducing resonance was to wrap the recording horn with tape.
Even drums, if planned and placed properly, could be effectively recorded and heard on even the earliest jazz and military band recordings. The loudest instruments such as the drums and trumpets were positioned the farthest away from the collecting horn. Lillian Hardin Armstrong, a member of King Oliver's Creole Jazz Band, which recorded at Gennett Records in 1923, remembered that at first Oliver and his young second trumpet, Louis Armstrong, stood next to each other and Oliver's horn could not be heard. "They put Louis about fifteen feet over in the corner, looking all sad."
Electrical recording
During the first half of the 1920s, engineers at Western Electric, as well as independent inventors such as Orlando Marsh, developed technology for capturing sound with a microphone, amplifying it with vacuum tubes (known as valves in the UK), and then using the amplified signal to drive an electromechanical recording head. Western Electric's innovations resulted in a broader and smoother frequency response, which produced a dramatically fuller, clearer and more natural-sounding recording. Soft or distant sounds that were previously impossible to record could now be captured. Volume was now limited only by the groove spacing on the record and the amplification of the playback device. Victor and Columbia licensed the new electrical system from Western Electric and recorded the first electrical discs during the spring of 1925. The first electrically recorded Victor Red Seal record was Chopin's "Impromptus" and Schubert's "Litanei" performed by pianist Alfred Cortot at Victor's studios in Camden, New Jersey.
A 1926 Wanamaker's ad in The New York Times offers records "by the latest Victor process of electrical recording". It was recognized as a breakthrough; in 1930, a Times music critic stated:
... the time has come for serious musical criticism to take account of performances of great music reproduced by means of the records. To claim that the records have succeeded in exact and complete reproduction of all details of symphonic or operatic performances ... would be extravagant ... [but] the article of today is so far in advance of the old machines as hardly to admit classification under the same name. Electrical recording and reproduction have combined to retain vitality and color in recitals by proxy.
The Orthophonic Victrola had an interior folded exponential horn, a sophisticated design informed by impedance-matching and transmission-line theory, and designed to provide a relatively flat frequency response. Victor's first public demonstration of the Orthophonic Victrola on 6 October 1925, at the Waldorf-Astoria Hotel was front-page news in The New York Times, which reported:
The audience broke into applause ... John Philip Sousa [said]: '[Gentlemen], that is a band. This is the first time I have ever heard music with any soul to it produced by a mechanical talking machine' ... The new instrument is a feat of mathematics and physics. It is not the result of innumerable experiments, but was worked out on paper in advance of being built in the laboratory ... The new machine has a range of from 100 to 5,000 [cycles per second], or five and a half octaves ... The 'phonograph tone' is eliminated by the new recording and reproducing process.
Sales of records plummeted precipitously during the early years of the Great Depression of the 1930s, and the entire record industry in America nearly foundered. In 1932, RCA Victor introduced a basic, inexpensive turntable called the Duo Jr., which was designed to be connected to their radio receivers. According to Edward Wallerstein (the general manager of the RCA Victor Division), this device was "instrumental in revitalizing the industry".
78 rpm materials
The production of shellac records continued throughout the 78 rpm era, which lasted until 1948 in industrialized nations.
During the Second World War, the United States Armed Forces produced thousands of 12-inch vinyl 78 rpm V-Discs for use by the troops overseas. After the war, the use of vinyl became more practical as new record players with lightweight crystal pickups and precision-ground styli made of sapphire or an exotic osmium alloy proliferated. In late 1945, RCA Victor began offering "De Luxe" transparent red vinylite pressings of some Red Seal classical 78s, at a de luxe price. Later, Decca Records introduced vinyl Deccalite 78s, while other record companies used various vinyl formulations trademarked as Metrolite, Merco Plastic, and Sav-o-flex, but these were mainly used to produce "unbreakable" children's records and special thin vinyl DJ pressings for shipment to radio stations.
78 rpm recording time
The playing time of a phonograph record is directly proportional to the available groove length divided by the turntable speed. Total groove length in turn depends on how closely the grooves are spaced, in addition to the record diameter. At the beginning of the 20th century, the early discs played for two minutes, the same as cylinder records. The 12-inch disc, introduced by Victor in 1903, increased the playing time to three and a half minutes. Because the standard 10-inch 78 rpm record could hold about three minutes of sound per side, most popular recordings were limited to that duration. For example, when King Oliver's Creole Jazz Band, including Louis Armstrong on his first recordings, recorded 13 sides at Gennett Records in Richmond, Indiana, in 1923, one side was 2:09 and four sides were 2:52–2:59.
In January 1938, Milt Gabler started recording for Commodore Records, and to allow for longer continuous performances, he recorded some 12-inch discs. Eddie Condon explained: "Gabler realized that a jam session needs room for development." The first two 12-inch recordings did not take advantage of their capability: "Carnegie Drag" was 3m 15s; "Carnegie Jump", 2m 41s. But at the second session, on 30 April, the two 12-inch recordings were longer: "Embraceable You" was 4m 05s; "Serenade to a Shylock", 4m 32s. Another way to overcome the time limitation was to issue a selection extending to both sides of a single record. Vaudeville stars Gallagher and Shean recorded "Mr. Gallagher and Mr. Shean", written by themselves or, allegedly, by Bryan Foy, as two sides of a 10-inch 78 in 1922 for Victor. Longer musical pieces were released as a set of records. In 1903 His Master's Voice in England made the first complete recording of an opera, Verdi's Ernani, on 40 single-sided discs.
In 1940, Commodore released Eddie Condon and his Band's recording of "A Good Man Is Hard to Find" in four parts, issued on both sides of two 12-inch 78s. The limited duration of recordings persisted from their advent until the introduction of the LP record in 1948. In popular music, the time limit of minutes on a 10-inch 78 rpm record meant that singers seldom recorded long pieces. One exception is Frank Sinatra's recording of Rodgers and Hammerstein's "Soliloquy", from Carousel, made on 28 May 1946. Because it ran 7m 57s, longer than both sides of a standard 78 rpm 10-inch record, it was released on Columbia's Masterwork label (the classical division) as two sides of a 12-inch record.
In the 78 era, classical-music and spoken-word items generally were released on the longer 12-inch 78s, about 4–5 minutes per side. For example, on 10 June 1924, four months after the 12 February premier of Rhapsody in Blue, George Gershwin recorded an abridged version of the seventeen-minute work with Paul Whiteman and His Orchestra. It was released on two sides of Victor 55225 and ran for 8m 59s.
Record albums
"Record albums" were originally booklets containing collections of multiple disc records of related material, the name being related to photograph albums or scrap albums. German record company Odeon pioneered the album in 1909 when it released the Nutcracker Suite by Tchaikovsky on four double-sided discs in a specially designed package. It was not until the LP era that an entire album of material could be included on a single record.
78 rpm releases in the microgroove era
In 1968, when the hit movie Thoroughly Modern Millie was inspiring revivals of Jazz Age music, Reprise planned to release a series of 78-rpm singles from their artists on their label at the time, called the Reprise Speed Series. Only one disc actually saw release, Randy Newman's "I Think It's Going to Rain Today", a track from his self-titled debut album (with "The Beehive State" on the flipside). Reprise did not proceed further with the series due to a lack of sales for the single, and a lack of general interest in the concept.
In 1978, guitarist and vocalist Leon Redbone released a promotional 78-rpm single featuring two songs ("Alabama Jubilee" and "Please Don't Talk About Me When I'm Gone") from his Champagne Charlie album.
In the same vein of Tin Pan Alley revivals, R. Crumb & His Cheap Suit Serenaders issued a number of 78-rpm singles on their Blue Goose record label. The most familiar of these releases is probably R. Crumb & His Cheap Suit Serenaders' Party Record (1980, issued as a "Red Goose" record on a 12-inch single), with the double-entendre "My Girl's Pussy" on the "A" side and the X-rated "Christopher Columbus" on the "B" side.
In the 1990s Rhino Records issued a series of boxed sets of 78-rpm reissues of early rock and roll hits, intended for owners of vintage jukeboxes. The records were made of vinyl, however, and some of the earlier vintage 78-rpm jukeboxes and record players (the ones that were pre-war) were designed with heavy tone arms to play the hard slate-impregnated shellac records of their time. These vinyl Rhino 78s were softer and would be destroyed by old juke boxes and old record players, but play well on newer 78-capable turntables with modern lightweight tone arms and jewel needles.
As a special release for Record Store Day 2011, Capitol re-released The Beach Boys single "Good Vibrations" in the form of a 10-inch 78-rpm record (b/w "Heroes and Villains"). More recently, The Reverend Peyton's Big Damn Band has released their tribute to blues guitarist Charley Patton Peyton on Patton on both 12-inch LP and 10-inch 78s.
New sizes and materials after WWII
CBS Laboratories had long been at work for Columbia Records to develop a phonograph record that would hold at least 20 minutes per side.
Research began in 1939, was suspended during World War II, and then resumed in 1945. Columbia Records unveiled the LP at a press conference in the Waldorf-Astoria on 21 June 1948, in two formats: in diameter, matching that of 78 rpm singles, and in diameter.
Unwilling to accept and license Columbia's system, in February 1949, RCA Victor released the first 45 rpm single, 7 inches in diameter with a large center hole. The 45 rpm player included a changing mechanism that allowed multiple disks to be stacked, much as a conventional changer handled 78s. Also like 78s, the short playing time of a single 45 rpm side meant that long works, such as symphonies and operas, had to be released on multiple 45s instead of a single LP, but RCA Victor claimed that the new high-speed changer rendered side breaks so brief as to be inconsequential. Early 45 rpm records were made from either vinyl or polystyrene. They had a playing time of eight minutes.
At first the two systems were marketed in competition, in what was called "The War of the Speeds".
Speeds
Shellac era
The older 78 rpm format continued to be mass-produced alongside the newer formats using new materials in decreasing numbers until the summer of 1958 in the U.S., and in a few countries, such as the Philippines and India (both countries issued recordings by the Beatles on 78s), into the late 1960s. For example, Columbia Records' last reissue of Frank Sinatra songs on 78 rpm records was an album called Young at Heart, issued in November 1954.
Microgroove and vinyl era
Columbia and RCA Victor each pursued their R&D secretly.
The commercial rivalry between RCA Victor and Columbia Records led to RCA Victor's introduction of what it had intended to be a competing vinyl format, the 7-inch (175 mm) 45 rpm disc, with a much larger center hole. For a two-year period from 1948 to 1950, record companies and consumers faced uncertainty over which of these formats would ultimately prevail in what was known as the "War of the Speeds" (see also Format war). In 1949 Capitol and Decca adopted the new LP format and RCA Victor gave in and issued its first LP in January 1950. The 45 rpm size was gaining in popularity, too, and Columbia issued its first 45s in February 1951. By 1954, 200 million 45s had been sold.
Eventually the 12-inch (300 mm) rpm LP prevailed as the dominant format for musical albums, and 10-inch LPs were no longer issued. The last Columbia Records reissue of any Frank Sinatra songs on a 10-inch LP record was an album called Hall of Fame, CL 2600, issued on 26 October 1956, containing six songs, one each by Tony Bennett, Rosemary Clooney, Johnnie Ray, Frank Sinatra, Doris Day, and Frankie Laine.
The 45 rpm discs also came in a variety known as extended play (EP), which achieved up to 10–15 minutes play at the expense of attenuating (and possibly compressing) the sound to reduce the width required by the groove. EP discs were cheaper to produce and were used in cases where unit sales were likely to be more limited or to reissue LP albums on the smaller format for those people who had only 45 rpm players. LP albums could be purchased one EP at a time, with four items per EP, or in a boxed set with three EPs or twelve items. The large center hole on 45s allows easier handling by jukebox mechanisms. EPs were generally discontinued by the late 1950s in the U.S. as three- and four-speed record players replaced the individual 45 players. One indication of the decline of the 45 rpm EP is that the last Columbia Records reissue of Frank Sinatra songs on 45 rpm EP records, called Frank Sinatra (Columbia B-2641) was issued on 7 December 1959.
The Seeburg Corporation introduced the Seeburg Background Music System in 1959, using a rpm 9-inch record with 2-inch center hole. Each record held 40 minutes of music per side, recorded at 420 grooves per inch.
From the mid-1950s through the 1960s, in the U.S. the common home record player or "stereo" (after the introduction of stereo recording) would typically have had these features: a three- or four-speed player (78, 45, , and sometimes rpm); with changer, a tall spindle that would hold several records and automatically drop a new record on top of the previous one when it had finished playing, a combination cartridge with both 78 and microgroove styli and a way to flip between the two; and some kind of adapter for playing the 45s with their larger center hole. The adapter could be a small solid circle that fit onto the bottom of the spindle (meaning only one 45 could be played at a time) or a larger adapter that fit over the entire spindle, permitting a stack of 45s to be played.
RCA Victor 45s were also adapted to the smaller spindle of an LP player with a plastic snap-in insert known as a "45 rpm adapter". These inserts were commissioned by RCA president David Sarnoff and were invented by Thomas Hutchison.
Capacitance Electronic Discs were videodiscs invented by RCA, based on mechanically tracked ultra-microgrooves (9541 grooves/inch) on a 12-inch conductive vinyl disc.
High fidelity
The term "high fidelity" was coined in the 1920s by some manufacturers of radio receivers and phonographs to differentiate their better-sounding products claimed as providing "perfect" sound reproduction. The term began to be used by some audio engineers and consumers through the 1930s and 1940s. After 1949 a variety of improvements in recording and playback technologies, especially stereo recordings, which became widely available in 1958, gave a boost to the "hi-fi" classification of products, leading to sales of individual components for the home such as amplifiers, loudspeakers, phonographs, and tape players. High Fidelity and Audio were two magazines that hi-fi consumers and engineers could read for reviews of playback equipment and recordings.
Stereophonic sound
A stereophonic phonograph provides two channels of audio, one left and one right. This is achieved by adding another vertical dimension of movement to the needle in addition to the horizontal one. As a result, the needle now moves not only left and right, but also up and down. But since those two dimensions do not have the same sensitivity to vibration, the difference needs to be evened out by having each channel take half its information from each direction by turning the channels 45 degrees from horizontal.
As a result of the 45-degree turn and some vector addition, it can be demonstrated that out of the new horizontal and vertical directions, one would represent the sum of the two channels, and the other representing the difference. Record makers decide to pick the directions such that the traditional horizontal direction codes for the sum. As a result, an ordinary mono disk is decoded correctly as "no difference between channels", and an ordinary mono player would simply play the sum of a stereophonic record without too much loss of information.
In 1957 the first commercial stereo two-channel records were issued first by Audio Fidelity followed by a translucent blue vinyl on Bel Canto Records, the first of which was a multi-colored-vinyl sampler featuring A Stereo Tour of Los Angeles narrated by Jack Wagner on one side, and a collection of tracks from various Bel Canto albums on the back.
Noise reduction systems
A similar scheme aiming at the high-end audiophile market, and achieving a noise reduction of about 20 to 25 dB(A), was the Telefunken/Nakamichi High-Com II noise reduction system being adapted to vinyl in 1979. A decoder was commercially available but only one demo record is known to have been produced in this format.
The availability of encoded disks in any of these formats stopped in the mid-1980s.
Yet another noise reduction system for vinyl records was the UC compander system developed by (ZWT) of (RFT). The system deliberately reduced disk noise by 10 to 12 dB(A) only to remain virtually free of recognizable acoustical artifacts even when records were played back without an UC expander. In fact, the system was undocumented yet introduced into the market by several East-German record labels since 1983. Over 500 UC-encoded titles were produced without an expander becoming available to the public. The only UC expander was built into a turntable manufactured by .
Formats
Types of records
The usual diameters of the holes on an EP record are .
Sizes of records in the United States and the UK are generally measured in inches, e.g. 7-inch records, which are generally 45 rpm records. LPs were 10-inch records at first, but soon the 12-inch size became by far the most common. Generally, 78s were 10-inch, but 12-inch and 7-inch and even smaller were made—the so-called "little wonders".
Standard formats
| Technology | Media and communication: Basics | null |
172161 | https://en.wikipedia.org/wiki/Blackcurrant | Blackcurrant | The blackcurrant (Ribes nigrum), also known as black currant or cassis, is a deciduous shrub in the family Grossulariaceae grown for its edible berries. It is native to temperate parts of central and northern Europe and northern Asia, where it prefers damp fertile soils. It is widely cultivated both commercially and domestically.
It is winter hardy, but cold weather at flowering time during the spring may reduce the size of the crop. Bunches of small, glossy black fruit develop along the stems in the summer and can be harvested by hand or by machine.
The raw fruit is particularly rich in vitamin C and polyphenols. Blackcurrants can be eaten raw but are usually cooked in sweet or savoury dishes. They are used to make jams, preserves, and syrups and are grown commercially for the juice market. The fruit is also used to make alcoholic beverages and dyes.
Description
Ribes nigrum is a medium-sized shrub, growing to . The leaves are alternate, simple, broad and long with five palmate lobes and a serrated margin. All parts of the plant are strongly aromatic. The flowers are produced in racemes known as "strigs" up to long containing 10–20 flowers, each about in diameter. Each flower has a hairy calyx with yellow glands, the five lobes of which are longer than the inconspicuous petals. There are five stamens surrounding the stigma and style and two fused carpels. The flowers open in succession from the base of the string and are mostly insect pollinated, but some pollen is distributed by the wind. A pollen grain landing on a stigma will germinate and send a slender pollen tube down the style to the ovule. In warm weather, this takes about 48 hours, but in cold weather, it may take a week, and by that time, the ovule may have passed the stage where it is receptive. If fewer than about 35 ovules are fertilised, the fruit may not be able to develop and will fall prematurely. Frost can damage both unopened and open flowers when the temperature falls below . The flowers at the base of the strig are more protected by the foliage and are less likely to be damaged.
In midsummer the strigs of green fruit ripen to edible berries, very dark purple in colour, almost black, with glossy skins and calyxes at the apex (the calyxes being persistent), each containing many seeds. An established bush can produce about of fruit each year.
Plants from Northern Asia are sometimes distinguished as a separate variety, , of which R. cyathiforme is considered a synonym.
Phytochemicals
Polyphenol phytochemicals present in the fruit, seeds, and leaves are being investigated for their potential biological activities. Major anthocyanins in blackcurrant pomace delphinidin-3-O-glucoside, delphinidin-3-O-rutinoside, cyanidin-3-O-glucoside, and cyanidin-3-O-rutinoside, which are retained in the juice concentrate are among other polyphenols.
Distribution and habitat
The blackcurrant is native to northern Europe and Asia.
Cultivation
Cultivation in Europe is thought to have started around the last decades of the 17th century.
Site selection and planting
Blackcurrants can grow well on sandy or heavy loams, or forest soils, as long as their nutrient requirements are met. They prefer damp, fertile but not waterlogged ground and are intolerant of drought. Although the bushes are winter hardy, frosts during the flowering period may adversely affect the yield and cold winds may restrict the number of flying insects visiting and pollinating the flowers. A soil pH of about 6 is ideal for blackcurrants and the ground can be limed if the soil is too acidic. Planting is usually done in the autumn or winter to allow the plants to become established before growth starts in the spring, but container-grown stock can be planted at any time of year.
Two-year-old bushes are usually planted but strong one-year-old stock can also be used. Planting certified stock avoids the risk of introducing viruses. On a garden scale the plants can be set at intervals of or they can be set in rows with planting intervals of and row separations of or more. In the UK, young bushes are generally planted deeper than their initial growing level to encourage new stems to grow from the base.
Manures and fertilizers
The blackcurrant requires a number of essential nutrients to thrive; nitrogen provides strong plant growth and stimulates the production of flower sprigs; phosphorus aids growth, the setting of fruit and crop yield; potassium promotes growth of individual shoots and increases the weight of individual fruits; magnesium is a constituent of chlorophyll and helps increase yields through interaction with potassium; calcium is required for cell division and enlargement and is particularly important for young plants and buds.
An annual spring mulch of well rotted manure is ideal and poultry manure can also be used but needs prior composting with straw or other waste vegetable material. Spent mushroom compost can be used but care should be taken as it often contains lime and blackcurrants prefer slightly acidic soils. The blackcurrant is a gross feeder and benefits from additional nitrogen, and phosphatic and potash fertilisers should also be applied annually. A balanced artificial fertilizer can be used and a 10-10-10 granular product can be spread around the bushes at the rate of per plant. Weed growth can be suppressed with an organic mulch such as sawdust, bark, mushroom compost or straw, heavy plastic topped with an organic mulch cover or landscape fabric.
Pruning
Blackcurrant fruit is borne primarily on one-year-old shoots. Newly planted bushes should be pruned severely, cutting all shoots back to two buds above ground level. This gives the plant a chance to get properly established before needing to put its energy into producing fruit. The general rule when pruning is to remove all weak shoots and those growing out sideways which may get weighed down when fruiting. The remaining branches should be thinned to remove old unproductive wood and to encourage new shoots. An established bush should not be allowed to become overcrowded and should have about one third of its main branches or stems removed each year. When harvesting by machine, plants with an upright growth habit are encouraged.
Harvesting
On a garden scale, the berries should be picked when dry and ripe. Commercially, most harvesting is done mechanically by straddle harvesters. These move continually down the rows, straddling a row of bushes, shaking the branches and stripping off the fruit. The blackcurrants are placed into half tonne bins and to minimise stoppage time, some machines have cross conveyors which direct the fruit into continuously moving trailers in the adjoining row. A modern machine can pick up to fifty tonnes of blackcurrants in a day using only one operator and two tractor drivers. The bins should be stored in a cool place. Some fruit is still picked by hand for use in the fresh fruit market.
Diseases and pests
Ribes plants are susceptible to several diseases and a number of insect pests. However, new varieties have been or are being developed to overcome some of these problems.
Reversion is a serious disease transmitted by the blackcurrant gall mite (Cecidophyopsis ribis). It causes a decline in yield and is quite widespread in Europe but is rarely encountered on other continents. Symptoms include a modification of leaf shape in summer and swollen buds ("big bud") in winter, each housing thousands of microscopic mites. As pest control has limited effectiveness, severely infected bushes should be destroyed. All new plants purchased should be certified as virus-free.
White pine blister rust (Cronartium ribicola) needs two alternate hosts to complete its life cycle. One host is plants in the genus Ribes. On the blackcurrant, it causes the leaves to become pale and later develop tiny orange pustules and sometimes a yellow filamentous coating on some leaves. The fruit crop is little affected but the leaves fall early and growth is slowed the following year. The other host is any of the white pines, in which it causes serious disease and mortality for the North American species that have not co-evolved with the rust. As a result, the blackcurrant was banned in the United States as a disease vector for much of the 20th century, and even after the federal ban was lifted in 1966, several U.S. states continued their own bans, some of which remain in force as of August 2021. The effectiveness of these restrictions is questionable, since other Ribes species also host the disease, some are native to North America, and others such as red currants and Ribes uva-crispa were never banned.
American gooseberry mildew and powdery mildew can infect the leaves and shoot tips, and botrytis may cause the fruit to rot in a wet season. Currant and gooseberry leaf spot (Drepanopeziza ribis) is another disease of blackcurrants, but it is not usually a serious problem as most cultivars now have some resistance.
The blackcurrant leaf midge can cause browning, crimping and distortion of leaves at the tips of shoots but it is seldom a serious problem. The blackcurrant sawfly (Nematus ribesii) lays its eggs on the underside of the leaves and the voracious larvae work their way along the shoots, stripping off leaf after leaf. In a serious attack, the bush can be denuded of leaves. Larvae of the currant borer drill their way along the centres of shoots, which wilt and die back. Other insect pests include scale insects, aphids and earwigs.
Research and cultivars
There are many cultivars of blackcurrant. 'Baldwin' was the mainstay of the industry for many years but it has now largely been superseded by more productive and disease-resistant varieties. During the 20th century in Europe, much hybridisation work has been carried out in order to reduce the plant's susceptibility to disease and frost and also to increase yields. This effort centered mainly in Scotland, Poland, and New Zealand.
In Britain the Scottish Crop Research Institute was tasked with developing new varieties suitable for growing in the north of the country. They produced new cultivars that had greater cold tolerance, especially in the spring, ripened earlier and more evenly and had greater fungal disease resistance. Frost tolerance was improved by selecting for late flowering and genetic research identified genes involved in resistance to gall mite and the blackcurrant reversion virus. 'Ben Lomond' was the first of the 'Ben' varieties and was released in 1975. This was followed by several other cultivars for the juicing industry such as 'Ben Alder' and 'Ben Tirran'. The cultivar 'Ben Hope' was released in 1998 with increased tolerance to gall mite, and in the same year, 'Ben Gairn' became available. It shows resistance to the reversion virus. For gardeners and the pick-your-own market, 'Ben Sarek', 'Ben Connan' and 'Big Ben' were introduced and have large, sweet berries. The cultivars 'Ben Connan' and 'Big Ben' have gained the Royal Horticultural Society's Award of Garden Merit. and new varieties are being developed continually to improve frost tolerance, disease resistance, machine harvesting, fruit quality, nutritional content and fruit flavour.
Varieties producing green fruit, less strongly flavoured and sweeter than typical blackcurrants, are cultivated in Finland, where they are called "greencurrants" (viherherukka). In Poland, the Research Institute of Horticulture has done work on improving the blackcurrant with regard to disease and pest resistance, fruit quality, adaptations to local conditions and mechanical harvesting. Researchers have crossed various varieties and introduced inter-specific genetic material from the gooseberry (Ribes grossularia), the redcurrant (Ribes rubrum) and the flowering currant (Ribes sanguineum). The resulting offspring were further back-crossed to R. nigrum. Cultivars produced include 'Tisel' and 'Tiben' in 2000 and 'Ores', 'Ruben' and 'Tines' in 2005. Further cultivars 'Polares' and 'Tihope' are being tested. Since 1991, New Zealand has become an important centre for research and development, as its temperate climate is particularly suitable for cultivation of the crop. Breeding programmes are concentrating on yield, large fruit size, consistency of cropping and upright habit.
In North America, there is a need for this fruit to have resistance to white pine blister rust. New cultivars such as 'Crusader', 'Coronet' and 'Consort' have been developed there by crossing R. nigrum with R. ussuriense and these show resistance to the disease. However the quality and yield of these varieties are poor as compared to non-resistant strains and only Consort is reliably self-fertile. Back-crossing these varieties to a parent have produced new strains such as 'Titania' that have a higher yield, better disease resistance, are more tolerant of adverse weather conditions and are suitable for machine harvesting. Two new releases from a black currant breeding program in British Columbia, Canada, 'Blackcomb' and 'Tahsis', were selected for their immunity to white pine blister rust and their frost tolerance.
Uses
Nutrition
Raw blackcurrants are 82% water, 15% carbohydrates, 1% protein and 0.4% fat (table). Per 100 g serving providing 63 kilocalories, the raw fruit has high vitamin C content (218% of the Daily Value, DV) and moderate levels of iron and manganese (12% DV each). Other nutrients are present in negligible amounts (less than 10% DV, table).
Blackcurrant seed oil is rich in vitamin E and unsaturated fatty acids, including alpha-linolenic acid and gamma-linolenic acid.
History
Decoction of the leaves, bark or roots was used as a traditional remedy.
During World War II, most fruits rich in vitamin C, such as oranges, became difficult to obtain in the United Kingdom. Since blackcurrant berries are a rich source of the vitamin, and blackcurrant plants are suitable for growing in the UK climate, the British Government encouraged their cultivation and soon the yield of the nation's crop increased significantly. From 1942 onwards, blackcurrant syrup was distributed free of charge to children under the age of two. This may have given rise to the lasting popularity of blackcurrant as a flavouring in Britain. In Britain the commercial crop is completely mechanised and about 1,400 hectares of the fruit are grown, mostly under contract to the juicing industry. Commercially, most large-scale cultivation of blackcurrants is done in eastern Europe for the juice and juice concentrate market. , major cultivation efforts to improve fruit characteristics occurred in Scotland, New Zealand, and Poland.
Blackcurrants were once popular in the United States as well, but became less common in the 20th century after currant farming was banned in the early 1900s, when blackcurrants, as a vector of white pine blister rust, were considered a threat to the U.S. logging industry. The federal ban on growing currants was shifted to the jurisdictions of individual states in 1966, and was lifted in New York State in 2003 through the efforts of horticulturist Greg Quinn. As a result, currant growing is making a comeback in New York, Vermont, Connecticut, California, and Oregon. However, several statewide bans still exist . Since the American federal ban curtailed currant production nationally for nearly a century, the fruit remains largely unknown in the United States and has yet to regain its previous popularity to levels enjoyed in Europe or New Zealand. Owing to its unique flavour and richness in polyphenols, dietary fibre and essential nutrients, awareness and popularity of blackcurrant is once again growing, with a number of consumer products entering the U.S. market.
Culinary
The fruit of blackcurrants when eaten raw has a strong, tart flavour. It can be made into jams and jellies which set readily because of the fruit's high content of pectin and acid. For culinary use, the fruit is usually cooked with sugar to produce a purée, which can then be passed through muslin to separate the juice. The purée can be used to make blackcurrant preserves and be included in cheesecakes, yogurt, ice cream, desserts, sorbets, and many other sweet dishes. The exceptionally strong flavour can be moderated by combining it with other fruits, such as raspberries and strawberries in summer pudding, or apples in crumbles and pies. The juice can be used in syrups and cordials. Blackcurrants are a common ingredient of rødgrød, a popular kissel-like dessert in North German and Danish cuisines.
Blackcurrants are also used in savoury cooking. Their astringency creates added flavour in sauces, meats and other dishes. Blackcurrants are included in some unusual combinations of foods. They can be added to tomato and mint to make a salad. Blackcurrants may accompany roast beef, grilled lamb, duck, seafood and shellfish. Canvasback duck with blackcurrants was a delicacy in nineteenth century New York. They can provide a dipping sauce at barbecues. They can be blended with mayonnaise, and used to invigorate bananas and other tropical fruits. Blackcurrants can be combined with dark chocolate or added to mincemeat in traditional mince pies at Christmas.
Japan imports US$3.6 million of New Zealand blackcurrants for uses as dietary supplements, snacks, functional food products and as quick-frozen (IQF) produce for culinary production as jams, jellies or preserves.
Beverages
The juice forms the basis for various squashes, juice drinks, and smoothies. In Britain, 95% of the blackcurrants grown are used to manufacture Ribena (a brand of fruit juice whose name is derived from Ribes nigrum) and similar fruit syrups and juices. Macerated blackcurrants are also the primary ingredient in the apéritif, crème de cassis, which in turn is added to white wine to produce a Kir or to champagne to make a Kir Royal.
In the UK, a blackcurrant squash may be mixed with beer or alcoholic cider to make drinks including "cider and black", "lager and black", or "snakebite and black".
In Russia, blackcurrant leaves may be used for flavoring tea or preserves, such as salted cucumbers, and berries for home winemaking. Sweetened vodka may also be infused with blackcurrant leaves making a deep greenish-yellow beverage with a tart flavor and astringent taste. The berries may be infused in a similar manner.
In the Netherlands, blackcurrants are used in a carbonated soft drink named "cassis", not to be confused with the alcoholic crème de cassis liqueur. The variety by Hero has been made since 1938 with blackcurrant juice concentrate as well as a small quantity of fermented blackcurrant juice.
Blackcurrant seed oil is an ingredient in cosmetics preparations, often in combination with vitamin E. The leaves can be extracted to yield a yellow dye, and the fruit is a source for a blue or violet dye resulting from its rich content of anthocyanins.
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172174 | https://en.wikipedia.org/wiki/R%C3%A9seau%20Express%20R%C3%A9gional | Réseau Express Régional | The Réseau Express Régional (; ), commonly abbreviated RER (), is a hybrid commuter rail and rapid transit system, similar to the S-Bahns of German-speaking countries and the S Lines of Milan, serving Paris and its suburbs. It acts as a combined city-center underground rail system and suburbs-to-city-center commuter rail. In the city center, it acts as a faster counterpart of the Paris Métro, having fewer stops.
Conceived of as a "métropolitain express" (express metro) during the mid 1930s, the scheme was revived in the 1950s and construction began in the early 1960s. The RER was not fully conceptualised until the completion of the Schéma directeur d'aménagement et d'urbanisme (roughly: "master plan for urban development") in 1965. The RER network, which initially comprised two lines, was formally inaugurated on 8 December 1977 in a ceremony that was attended by President Valéry Giscard d'Estaing. A second phase of construction commenced at the end of the 1970s which saw additional lines constructed along with extensions to the original two. The RER is operated partly by RATP, the authority that operates most of the public transport in Paris, and partly by the SNCF, France's national rail operator.
As of 2023, the network consists of five lines: A, B, C, D and E. The network has 257 stations and has interchanges with the Métro and commuter rail within the City of Paris and the suburbs. The lines are identified by letters to avoid confusion with the Métro lines, which are identified by numbers. The network is still expanding: RER E, which opened in 1999, is planned for westward extension toward La Défense and Mantes-la-Jolie in two phases by 2024–2026. The performance of the RER has made it a model for proposals to improve transit within other cities. In November 2022, French President Emmanuel Macron announced the creation of the additional RERs that will serve the ten largest French metropolises by 2040.
Characteristics
The RER contains 257 stations, 33 of which are within the city of Paris, and runs over of track, including underground. Each line passes through the city almost wholly underground and on tracks dedicated to the RER, but some city center tracks are shared between line D and line B. The RER is operated partly by RATP, the authority that operates most public transport in Paris, and partly by SNCF, the national rail operator. The system, which is structured in a traditional radial arrangement, operates a through-service and uses a single fare model that works seamlessly with several other public transit systems. Total traffic on the central sections of lines A and B, operated by RATP, was 452 million people in 2006; in the same year, total traffic on all Paris area commuter lines operated by SNCF (both RER and Transilien trains) was 657 million.
RATP manages 65 RER stations, including all stations on Line A east of Nanterre-Préfecture and those on the branch to Saint-Germain-en-Laye. It also operates stations on Line B south of Gare du Nord. Other stations on the two lines and those on lines C, D and E are operated by SNCF. Of the RER stations operated by RATP, 9 have interchanges with Métro lines, and 9 allow transfer to SNCF's Transilien service. In comparison to the Metro, the RER provides better coverage of Paris' suburbs and typically operates at higher speeds and with greater distances between stations. Within the city center, RER services practice limited stop operations.
History
Origins
Its roots are in the 1936 Ruhlmann-Langewin plan of the Compagnie du chemin de fer métropolitain de Paris (Metropolitan Railway Company of Paris) for a "métropolitain express" (express metro). The company's post-war successor, RATP, revived the scheme in the 1950s, and in 1960 an interministerial committee decided to go ahead with the construction of an east-west line. Subsequently, the central part of the RER was completed between 1962 and 1977 in a large-scale civil engineering project whose chief supervisor was Siavash Teimouri. The construction of the RER was a major undertaking, being highly visible to both Parisians and visiting tourists at various sites across the city for an extended period.
As its instigator, RATP was granted authority to run the new link. The embryonic (and as yet unnamed) RER was not properly conceived until the 1965 Schéma directeur d'aménagement et d'urbanisme (roughly: "master plan for urban development"), which envisioned an H-shaped network with two north-south routes. Between 1969 and 1970, RATP purchased the Vincennes and Saint-Germain lines from SNCF, as the basis for the east-west link. Only a single north-south route crossing the Left Bank has so far come to fruition, although the Métro's line 13 has been extended to perform a similar function.
The RER's first phase of construction during the 1960s and 1970s was marked by scale and expense. In 1973 alone, FRF 2 billion were committed to the project in the budget, equating to roughly €1.37 billion in 2005 terms, and closer to double that as a proportion of the region's (then much smaller) economic output. The construction cost was controversial at the time of its construction. This initial investment, along with subsequent spending, was partly financed by the versement transport, a local tax levied on businesses that was introduced in July 1971. It has remained in effect into the twenty-first century.
First phase
During the first phase of construction, the Vincennes and Saint-Germain lines became the ends of the east-west Line A, the central section of which was opened station by station between 1969 and 1977. On its completion, Line A was joined by the initial southern section of the north-south Line B. Both Lines A and B were intentionally designed to converge with as many of the existing commuter lines as possible as to maximise its usefulness to existing travelers. During this first phase, six new stations were built, three of which are entirely underground.
Construction was inaugurated by Robert Buron, then Minister for Public Works, at the Pont de Neuilly on 6 July 1961, four years before the publication of the official network blueprint. The rapid expansion of the La Défense business district in the west made the western section of the first line a priority. Nation, the first new station, was opened on 12 December 1969 and became for the next 8 years the new western terminus of the Vincennes line. The section from Étoile (not yet renamed after Charles de Gaulle) to was opened a few weeks later. It was later extended eastward to the newly built on 23 November 1971, and westward to Saint-Germain-en-Laye on 1 October 1972. The latter extension was achieved by a connection to the existing Saint-Germain-en-Laye line, the oldest railway line into Paris, at Nanterre.
The RER network was inaugurated on 8 December 1977 with the joining of the eastern -Boissy segment and the western - segment at . The inauguration was attended by President Valery Giscard D'Estaing. The southern Ligne de Sceaux was simultaneously extended from Luxembourg to meet Line A at Châtelet – Les Halles, becoming the new Line B. The system of line letters was introduced to the public on this occasion, though it had been used internally by RATP and SNCF for some time.
Second phase
A second phase of construction commenced at the end of the 1970s, which was carried out at a slower pace than the first phase. SNCF gained the authorisation to operate its own routes, which became lines C, D and E. Extensive sections of suburban tracks were added to the network, but only four new stations were built.
During 1979, Line C (along the Left Bank of the Seine) was added, although it almost entirely comprised existing SNCF lines. The main civil engineering works performed involved the construction of a connecting link between Invalides and Musée d'Orsay. In 1981, Line B was extended through to Gare du Nord via a new deep tunnel from Châtelet – Les Halles. It was later extended further northward. By 1992, a total of 233 miles of track was operational, while a further 94 miles were under construction.
During 1995, Line D (north to south-east via Châtelet – Les Halles) was completed; its primary feature was a purpose-built deep tunnel between Châtelet – Les Halles and Gare de Lyon. No new building work was necessary at Châtelet – Les Halles, as additional platforms for Line D had been built at the time of the station's construction 20 years earlier. In 1999, Line E was added to the network, connecting the north-east with Gare Saint-Lazare by means of a new deep tunnel from Gare de l'Est.
Maps
Rolling stock
The predominance of suburban SNCF track on the RER network explains why RER trains use overhead line power and run on the left, like SNCF trains (except in Alsace-Moselle), contrary to the Métro where trains use third rail power and run on the right. RER trains run by the two different operators share the same track infrastructure, a practice called interconnection. On the RER, interconnection required the development of specific trains (MI 79 series for Materiel d'Interconnexion 1979, and MI 2N series for Materiel d'Interconnexion à 2 niveaux (double-deck interconnection stock)) capable of operating under both 1.5 kV direct current on the RATP network and 25 kV / 50 Hz alternating current on the SNCF network. The MS 61 series (Matériel Suburbain 1961) can be used only on the 1.5 kV DC network.
The RER's tunnels have unusually large cross-sections. This is due to a 1961 decision to build them according to a loading gauge standard created by the Union Internationale des Chemins de Fer (UIC), with space for overhead catenary power supply to trains. Single-track tunnels measure 6.30 m across and double-track tunnels up to 8.70 m, meaning a cross-sectional area of up to 50 square metres, larger than that of the stations on many comparable underground rail networks.
The first RER rolling stock in fact predated the formation of the RER by 40 years, with the Z 23000 stock used on the ligne de Sceaux (which was thereafter integrated into RER B) from 1937 until 27 February 1987. In 1965 the Z 5300 train was introduced, followed by the MS 61 in 1967, MI 79 in 1980, MI 84 and Z 8800 in 1985, Z 20500 in 1988, MI 2N in 1996, Z 20900 in 2001, MI 09 on 5 December 2011, Z 50000 (Francilien) in 2015 and Regio 2N (Z 57000) in 2019. In 2017, it was announced that an consortium comprising Alstom and Bombardier Transportation had been selected to supply 255 X’Trapolis Cityduplex double-deck electric multiple units to replace aging rolling stock on both lines D and E under a €3.75 billion arrangement. In April 2023, 60 additional Cityduplexs were ordered to increase service frequency.
Many services are performed by double-decker train sets, usually operating in double formations.
Lines
Stations
Ten new stations have been built under the heart of Paris since the 1960s as part of the RER project. The six stations of Line A opened between 1969 and 1977 are:
(1969): deep construction at the Place de la Nation.
(1970): deep construction at the Arc de Triomphe.
(1970): near-surface construction beneath the site of the yet-to-be-built Grande Arche de la Défense, just outside the Paris city boundary.
(1971): deep construction near Gare Saint-Lazare.
(1977): near-surface construction on the site of the former marketplace, claimed in 2017 to be the largest underground station in Europe.
Gare de Lyon (1977): near-surface construction beneath and alongside the main-line SNCF station.
Some controversy followed the construction of the Line A. Using the model of the existing Métro, and unlike any other underground network in the world, engineers elected to build the three new deep stations (Étoile, and ) as single monolithic halls with lateral platforms and no supporting pillars. A hybrid solution of adjacent halls was rejected on the grounds that it "completely sacrificed the architectural aspect" of the oeuvre. The scale in question was vast: the new stations cathédrales were up to three times longer, wider and taller than Métro stations, and hence 20 or 30 times more voluminous. Most importantly, unlike the Métro they were to be constructed deep underground. The decision turned out to be expensive: around 8 billion francs for the three stations, equivalent to €1.2 billion in 2005 terms, with the two-level the costliest of the three. The comparison was obvious and unfavourable with London's Victoria line, a deep line of constructed during the same period using a two-tunnel approach at vastly lower cost though with a lower capacity. However, the three stations represent undeniable engineering feats and are noticeably less claustrophobic than traditional underground stations.
Only two stations were inaugurated to complete Lines B, C and D:
Gare du Nord (1982): near-surface construction on two levels.
St-Michel - Notre-Dame (1988): deep construction on an existing stretch of the Line B between Luxembourg and Châtelet - Les Halles with two tunnels, common in other underground systems but unique in Paris. The station is situated on and built at the same time as the Luxembourg Châtelet tunnel.
Two stations were added to the network as part of Line E in the 1990s. They are notable for their lavishly spacious deep construction, a technique not used since . Although similar to the three 1960s "cathedral stations" of Line A, their passenger traffic has so far proved vastly lower.
(1999): deep construction serving both Gare du Nord and Gare de l'Est.
(1999): deep construction serving Gare Saint-Lazare and Auber.
Usage
Journey times, particularly on east-west and north-south routes, have been reduced by the RER; in combination with the cross-platform connection at Châtelet - Les Halles, even certain "diagonal" trips have reduced journey times. Typically being used for leisure journeys, the RER has made a significant social impact. By bringing far-flung suburbs within easy reach of central Paris, the network has aided the reintegration of the traditionally insular capital with its periphery. Evidence of this social impact can be seen at Châtelet - Les Halles, whose neighbourhood and Forum des Halles leisure and shopping facilities are popular among banlieusards (suburbanites), in particular from eastern suburbs.
Lines A and B reached saturation relatively quickly, exceeding by far all traffic expectations: up to 55,000 passengers per hour in each direction on Line A (1992), the highest such figure outside of East Asia. Despite a frequency of more than one train every two minutes, made possible by the installation of digital signalling in 1989, and the gradual introduction of double-decker trains from 1998 to 2017, the central stations of Line A are critically crowded at peak times. In June 2015, a contract valued at €20 million was awarded to Alstom Transport to develop and install an automatic train operation (ATO) system on RER A; at the time, this line was the most heavily frequented regional line in Europe, the introduction of ATO enabled increased service frequency and improved performance to be achieved.
The RER has a substantial impact on the suburban areas of greater Paris, specifically on land values near the stations along its lines. The RER has received criticism for its high level of particle pollution during busy periods, largely due to train braking. Pollution by PM10 particles regularly reaches 400 μg/m3 at , much more than at neighboring metro stations and eight times the EU Commission's daily average limit of 50 μg/m3.
| Technology | France | null |
172273 | https://en.wikipedia.org/wiki/Tick | Tick | Ticks are parasitic arachnids of the order Ixodida. They are part of the mite superorder Parasitiformes. Adult ticks are approximately 3 to 5 mm in length depending on age, sex, species, and "fullness". Ticks are external parasites, living by feeding on the blood of mammals, birds, and sometimes reptiles and amphibians. The timing of the origin of ticks is uncertain, though the oldest known tick fossils are from the Cretaceous period, around 100 million years old. Ticks are widely distributed around the world, especially in warm, humid climates.
Ticks belong to two major families, the Ixodidae or hard ticks, and the Argasidae, or soft ticks. Nuttalliella, a genus of tick from southern Africa, is the only member of the family Nuttalliellidae, and represents the most primitive living lineage of ticks. Adults have ovoid/pear-shaped bodies (idiosomas) which become engorged with blood when they feed, and eight legs. Their cephalothorax and abdomen are completely fused. In addition to having a hard shield on their dorsal surfaces, known as the scutum, hard ticks have a beak-like structure at the front containing the mouthparts, whereas soft ticks have their mouthparts on the underside of their bodies. Ticks locate potential hosts by sensing odor, body heat, moisture, and/or vibrations in the environment.
Ticks have four stages to their life cycle, namely egg, larva, nymph, and adult. Ticks belonging to the Ixodidae family undergo either a one-host, two-host, or three-host life cycle. Argasid ticks have up to seven nymphal stages (instars), each one requiring blood ingestion, and as such, Argasid ticks undergo a multihost life cycle. Because of their hematophagous (blood-ingesting) diets, ticks act as vectors of many serious diseases that affect humans and other animals.
Biology
Taxonomy and phylogeny
Ticks belong to the Parasitiformes, a distinctive group of mites that are separate from the main group of mites, the Acariformes. Whether the two groups are more closely related to each other than to other arachnids is uncertain, and studies often recover them as not closely related. Within the Parasitiformes, ticks are most closely related to the Holothyrida, a small group of free living scavengers with 32 described species confined to the landmasses that formed the supercontinent Gondwana.
Relationships among members of the Parasitiformes, after Klompen, 2010:
Fossilized ticks have been discovered from the end of the Early Cretaceous onwards, most commonly in amber. The oldest discovered tick fossils are an argasid bird tick from Late Cretaceous (Turonian ~94-90 million years ago) aged New Jersey amber, and various ticks found in Burmese amber, including Khimaira which does not belong to any living family of tick, the living genus Nuttalliella and the possible nuttalliellid genera Deinocroton and Legionaris, as well as the members of the living ixodid genera Amblyomma, Ixodes, Haemaphysalis, Bothriocroton and Archaeocroton dating the earliest Cenomanian stage of the Late Cretaceous, around . An undescribed juvenile tick is known from late Albian amber, dating to 105 million years ago. The younger Baltic and Dominican ambers have also yielded examples that can be placed in living genera. A phylogenetic analysis suggests that the last common ancestor of all living ticks likely lived around 195 million years ago in the Southern Hemisphere, in what was then Gondwana.
Ticks belong to three different families. The majority of tick species belong to the two families: Ixodidae (hard ticks) and Argasidae (soft ticks). The third living family is Nuttalliellidae, named for the bacteriologist George Nuttall. It comprises a single species, Nuttalliella namaqua, and as such is a monotypic taxon. Nuttalliella namaqua is found in southern Africa ranging from Tanzania to Namibia and South Africa.
Relationships of living and extinct tick families, after Chitimia-Dobler et al. 2022:The Ixodidae contain over 700 species of hard ticks with a scutum or hard shield, which the Argasidae lack. The Argasidae contain about 200 species; the genera accepted are Antricola, Argas, Nothoaspis, Ornithodoros, and Otobius. They have no scutum, and the capitulum (mouth and feeding parts) is concealed beneath the body. The phylogeny of the Ixodida within the Acari is shown in the cladogram, based on a 2014 maximum parsimony study of amino acid sequences of 12 mitochondrial proteins. The Argasidae appear monophyletic in this study.
Anatomy and physiology
Ticks, like mites, belong to the subclass Acari that lack their primary somatic segmentation of the abdomen (or opisthosoma), rather these parasitic arachnids present a subsequent fusion of the abdomen with the cephalothorax (or prosoma). The tagmata typical of other Chelicerata have developed into the gnathosoma (head), which is retractable and contains the mouthparts, and idiosoma (body), which contains the legs, digestive tract, and reproductive organs. The gnathosoma is a feeding structure with mouthparts adapted for piercing skin and sucking blood; it is the front of the head and contains neither the brain nor the eyes. Features of the gnathosoma include two palps, two chelicerae, and hypostome. The hypostome acts as stabilizer and helps to anchor the tick's mouthparts to the host. The chelicerae are specialized appendages used for cutting and piercing into the host's skin while palps are leglike appendages that are sensory in function.
The ventral side of the idiosoma bears sclerites, and the gonopore is located between the fourth pair of legs. In the absence of segmentation, the positioning of the eyes, limbs, and gonopore on the idiosoma provide the only locational guidance.
Larval ticks hatch with six legs, acquiring the other two after a blood meal and molting into the nymph stage. In the nymphal and adult stages, ticks have eight legs, each of which has seven segments and is tipped with a pair of claws. The legs are sometimes ornamented and usually bear sensory or tactile hairs. In addition to being used for locomotion, the tarsus of leg I contains a unique sensory structure, Haller's organ, which can detect odors and chemicals emanating from the host, as well as sensing changes in temperature and air currents. Ticks can also use Haller's organs to perceive infrared light emanating from a host. When stationary, their legs remain tightly folded against the body.
Ticks are extremely resilient animals. They can survive in a near vacuum for as long as half an hour. Their slow metabolism during their dormant periods enables them to go prolonged durations between meals. Even after 18 weeks of starvation, they can endure repeated two-day bouts of dehydration followed by rehydration, but their survivability against dehydration drops rapidly after 36 weeks of starvation. To keep from dehydrating, ticks hide in humid spots on the forest floor or absorb water from subsaturated air by secreting hygroscopic fluid produced by the salivary glands onto the external mouthparts and then reingesting the water-enriched fluid.
Ticks can withstand temperatures just above for more than two hours and can survive temperatures between for at least two weeks. Ticks have even been found in Antarctica, where they feed on penguins.
Most ticks are plain brown or reddish brown. However, the scuta of some species are decorated with white patterns.
Ixodidae
In nymphs and adults, the is prominent and projects forwards from the body. The eyes are close to the sides of the scutum and the large spiracles are located just behind the coxae of the fourth pair of legs. The hard protective scutellum, a characteristic of this family, covers nearly the whole dorsal surface in males, but is restricted to a small, shield-like structure behind the capitulum in females and nymphs. When an ixodid attaches to a host the bite is typically painless and generally goes unnoticed. They remain in place until they engorge and are ready to molt; this process may take days or weeks. Some species drop off the host to molt in a safe place, whereas others remain on the same host and only drop off once they are ready to lay their eggs.
Argasidae
The body of a soft tick is pear-shaped or oval with a rounded anterior portion. The mouthparts cannot be seen from above, as they are on the ventral surface. A centrally positioned dorsal plate with ridges projecting slightly above the surrounding surface, but with no decoration are often present. Soft ticks possess a leathery cuticle as well. A pattern of small, circular depressions expose where muscles are attached to the interior of the integument. The eyes are on the sides of the body, the spiracles open between legs 3 and 4, and males and females only differ in the structure of the genital pore.
Nuttalliellidae
Nuttalliellidae can be distinguished from both ixodid and argasid ticks by a combination of a projecting gnathosoma and a soft leathery skin. Other distinguishing characteristics include the position of the stigmata, the lack of setae, the strongly corrugated integument, and the form of the fenestrated plates.
Diet and feeding
Ticks are ectoparasites and consume blood to satisfy all of their nutritional requirements. They are obligate hematophages, and require blood to survive and move from one stage of life to another. Ticks can fast for long periods of time, but eventually die if unable to find a host. Hematophagy evolved independently at least six times in arthropods living during the late Cretaceous; in ticks it is thought to have evolved 120 million years ago through adaptation to blood-feeding. This behavior evolved independently within the separate tick families as well, with differing host-tick interactions driving the evolutionary change.
Some ticks attach to their host rapidly, while others wander around searching for thinner skin, such as that in the ears of mammals. Depending on the species and life stage, preparing to feed can take from ten minutes to two hours. On locating a suitable feeding spot, the tick grasps the host's skin and cuts into the surface. It extracts blood by cutting a hole in the host's epidermis, into which it inserts its hypostome and prevents the blood from clotting by excreting an anticoagulant or platelet aggregation inhibitor.
Ticks find their hosts by detecting an animals' breath and body odors, sensing body heat, moisture, or vibrations. A common misconception about ticks is they jump onto their host; however, they are incapable of jumping, although static electricity from their hosts has been shown to be capable of pulling the tick over distances several times their own body length. Many tick species, particularly Ixodidae, lie in wait in a position known as "questing". While questing, ticks cling to leaves and grasses by their third and fourth pairs of legs. They hold the first pair of legs outstretched, waiting to grasp and climb on to any passing host. Tick questing heights tend to be correlated with the size of the desired host; nymphs and small species tend to quest close to the ground, where they may encounter small mammalian or bird hosts; adults climb higher into the vegetation, where larger hosts may be encountered. Some species are hunters and lurk near places where hosts may rest. Upon receiving an olfactory stimulus or other environmental indication, they crawl or run across the intervening surface.
Other ticks, mainly the Argasidae, are nidicolous, finding hosts in their nests, burrows, or caves. They use the same stimuli as non-nidicolous species to identify hosts, with body heat and odors often being the main factors. Many of them feed primarily on birds, though some Ornithodoros species, for example, feed on small mammals. Both groups of soft tick feed rapidly, typically biting painfully and drinking their fill within minutes. Unlike the Ixodidae that have no fixed dwelling place except on the host, they live in sand, in crevices near animal dens or nests, or in human dwellings, where they come out nightly to attack roosting birds or emerge when they detect carbon dioxide in the breath of their hosts.
Ixodidae remain in place until they are completely engorged. Their weight may increase by 200 to 600 times compared to their prefeeding weight. To accommodate this expansion, cell division takes place to facilitate enlargement of the cuticle. In the Argasidae, the tick's cuticle stretches to accommodate the fluid ingested, but does not grow new cells, with the weight of the tick increasing five- to tenfold over the unfed state. The tick then drops off the host and typically remains in the nest or burrow until its host returns to provide its next meal.
Tick saliva contains about 1,500 to 3,000 proteins, depending on the tick species. The proteins with anti-inflammatory properties, called evasins, allow ticks to feed for eight to ten days without being perceived by the host animal. Researchers are studying these evasins with the goal of developing drugs to neutralise the chemokines that cause myocarditis, heart attack, and stroke.
Ticks do not use any other food source than vertebrate blood and therefore ingest high levels of protein, iron and salt, but few carbohydrates, lipids or vitamins. Ticks’ genomes have evolved large repertoires of genes related to this nutritional challenge, but they themselves cannot synthesize the essential vitamins that are lacking in blood meal. To overcome these nutritional deficiencies, ticks have evolved obligate interactions with nutritional endosymbionts. The first appearance of ticks and their later diversification were largely conditioned by this nutritional endosymbiosis lasting for millions of years. The most common of these nutritional endosymbionts belong to the Coxiella and Francisella bacterial genera. These intracellular symbiotic microorganisms are specifically associated with ticks and use transovarial transmission to ensure their persistence. Although Coxiella and Francisella endosymbionts are distantly related bacteria, they have converged towards an analogous B vitamin-based nutritional mutualism with ticks. Their experimental elimination typically results in decreased tick survival, molting, fecundity and egg viability, as well as in physical abnormalities, which all are fully restored with an oral supplement of B vitamins. The genome sequencing of Coxiella and Francisella endosymbionts confirmed that they consistently produce three B vitamin types, biotin (vitamin B7), riboflavin (B2) and folate (B9). As they are required for tick life cycle, these obligate endosymbionts are present in all individuals of the tick species they infect, at least at early stages of development since they may be secondarily lost in males during nymphal development. Since Coxiella and Francisella endosymbionts are closely related to pathogens, there is a substantial risk of misidentification between endosymbionts and pathogens, leading to an overestimation of infection risks associated with ticks.
Range and habitat
Tick species are widely distributed around the world. They tend to flourish more in warm, humid climates, because they require a certain amount of moisture in the air to undergo metamorphosis, and low temperatures inhibit their development of eggs to larvae.
The occurrence of ticks and tick-borne illnesses in humans is increasing. Tick populations are spreading into new areas, due in part to the warming temperatures of climate change.
Tick parasitism is widely distributed among host taxa, including marsupial and placental mammals, birds, reptiles (snakes, iguanas, and lizards), and amphibians. Ticks of domestic animals cause considerable harm to livestock through pathogenic transmission, causing anemia through blood loss, and damaging wool and hides. The Tropical Bont tick wreaks havoc on livestock and wildlife in Africa, the Caribbean, and several other countries through the spread of disease, specifically heartwater disease. The spinose ear tick has a worldwide distribution, the young feed inside the ears of cattle and various wildlife.
A habitat preferred by ticks is the interface where a lawn meets the forest, or more generally, the ecotone, which is unmaintained transitional edge habitat between woodlands and open areas. Therefore, one tick management strategy is to remove leaf litter, brush, and weeds at the edge of the woods. Ticks like shady, moist leaf litter with an overstory of trees or shrubs and, in the spring, they deposit their eggs into such places allowing larvae to emerge in the fall and crawl into low-lying vegetation. The 3 meter boundary closest to the lawn's edge are a tick migration zone, where 82% of tick nymphs in lawns are found.
Ecology
In general, ticks are found wherever their host species occur. Migrating birds carry ticks with them on through their migrations; a study of migratory birds passing through Egypt discovered more than half the bird species examined were carrying ticks. It was also observed the tick species varied depending on the season of migration, in this study it is spring and autumn migrations, this is thought to occur due to the seasonal periodicities of the different species.
For an ecosystem to support ticks, it must satisfy two requirements; the population density of host species in the area must be great enough and it must be humid enough for ticks to remain hydrated. Due to their role in transmitting Lyme disease, Ixodid ticks, particularly the North American I. scapularis, have been studied using geographic information systems to develop predictive models for ideal tick habitats. According to these studies, certain features of a given microclimate – such as sandy soil, hardwood trees, rivers, and the presence of deer – were determined to be good predictors of dense tick populations.
Mites and nematodes feed on ticks, which are also a minor nutritional resource for birds. More importantly, ticks act as a disease vector and behave as the primary hosts of many different pathogens such as spirochaetes. Ticks carry various debilitating diseases therefore, ticks may assist in controlling animal populations and preventing overgrazing.
Ticks can transmit an array of infectious diseases that affect humans and other animals. Ticks that carry zoonotic pathogens often tend to have a wide host range. The infective agents can be present not only in the adult tick, but also in the eggs produced plentifully by the females. Many tick species have extended their ranges as a result of the movements of people, domesticated pets, and livestock. With increasing participation in outdoor activities such as wilderness hikes, more people and their dogs may find themselves exposed to ticks.
Life cycle
All three tick families ticks have four life cycle stages: egg, larva, nymph, and adult.
Ixodidae
Ixodidae ticks have three different life cycles. Depending on the species, Ixodids can either possess a one-host life cycle, two-host life cycle, or three-host life cycle.
One-host ticks
In one-host ticks the tick remains on the host through the larval, nymphal, and adult stages, only to leave the host to lay eggs. Eggs laid in the environment hatch into larvae, which immediately seek out a host in which to attach and feed. Fed larvae molt into unfed nymphs that remain on the host. After engorging on the host's blood, the nymphs molt into sexually mature adults that remain on the host in order to feed and mate. Once a female is both fed and ready to lay eggs, only then does she leave the host in search of a suitable area to deposit her eggs. Ticks that follow this life cycle are called one-host ticks. The winter tick Dermacentor albipictus and the cattle tick Boophilus microplus are examples of one-host ticks.
Two-host ticks
The life cycle of a two-host tick often spans two years. During fall the pregnant female tick will drop off her second host and lay her eggs. The eggs hatch during winter, the following spring the larvae emerge and attach to their first host. Newly hatched larvae attach to a host in order to obtain a blood meal. They remain on the host then develop into nymphs. Once engorged, they drop off the host and find a safe area in the natural environment in which to molt into adults, this typically occurs during the winter. Both male and female adults seek out a host on which to attach, which may be the same body that served as host during their early development but is often a larger mammal. Once attached, they feed and mate. Gravid females drop from the host to oviposit in the environment. Ticks that complete their life cycle in this manner are called two-host ticks, like Hyalomma anatolicum excavatum.
Three-host ticks
Most ixodid ticks require three hosts, and their life cycles typically span three years. The female tick drops off its host, often in the fall, and lays thousands of eggs. The larvae hatch in the winter and emerge in the spring. When the larvae emerge, they attach and feed primarily on small mammals and birds. During the summer the larvae become engorged and drop off the first host to molt and become nymphs, this often occurs during the fall. The following spring the nymphs emerge and seek out another host, often a small rodent. The nymphs become engorged and drop off the host in the fall to molt and become adults. The following spring the adult ticks emerge and seek out a larger host, often a large mammal such as cattle or even humans. Females will mate on their third host. Female adults then engorge on blood and prepare to drop off to lay her eggs on the ground, while males feed very little and remain on the host in order to continue mating with other females.
Argasidae
Argasid ticks, unlike ixodid ticks, may go through up to seven nymphal stages (instars), requiring a meal of blood each time. Often, egg laying and mating occurs detached from the host in a safe environment. The eggs hatch and the larvae feed on a nearby host for anywhere from a few hours to several days, this depends on the species of tick. After they feed the larvae drop and molt into their first nymphal instars, then the nymph seeks out and feeds on its second host, often this is the same as the first host, within an hour. This process occurs repeatedly and until the last nymphal instar occurs, thus allowing the tick to molt into an adult. Once an adult these ticks feed rapidly and periodically their entire life cycle. In some species an adult female may lay eggs after each feeding. Their life cycles range from months to years. The adult female argasid tick can lay a few hundred to over a thousand eggs over the course of her lifetime. Both male and female adults feed on blood, and they mate off the host. During feeding, any excess fluid is excreted by the coxal glands, a process that is unique to argasid ticks.
Nuttalliellidae
Nuttalliellidae is an elusive monotypic family of tick, that is, possesses a single species, Nuttalliella namaqua. There is little to nothing known about the life cycle and feeding habits of N. namaqua but it is speculated this species of tick has multiple different hosts.
Relationship with humans
Tick-borne disease
Ticks can transmit many kinds of pathogens, such as bacteria, viruses, and protozoa, that infect ticks’ hosts. A tick can harbor more than one type of pathogen, making diagnosis more difficult. Species of the bacterial genus Rickettsia are responsible for typhus, rickettsialpox, boutonneuse fever, African tick bite fever, Rocky Mountain spotted fever, Flinders Island spotted fever, and Queensland tick typhus (Australian tick typhus). Other tick-borne diseases include Lyme disease and Q fever, Colorado tick fever, Crimean–Congo hemorrhagic fever, tularemia, tick-borne relapsing fever, babesiosis, ehrlichiosis, Bourbon virus, and tick-borne meningoencephalitis, as well as bovine anaplasmosis and the Heartland virus. In the United States, Lyme disease is the most commonly reported vector-borne disease in the country.
Some species, notably the Australian paralysis tick, are also intrinsically venomous and can cause tick paralysis. Eggs can become infected with pathogens inside a female tick's ovaries, in which case the larval ticks are infectious immediately at hatching, before feeding on their first host. Tropical bont ticks transmit the heartwater, which can be particularly devastating in cattle. The ticks carried by migratory birds act as reservoirs and vectors of foreign infectious diseases. In the Egyptian migratory bird study, over 20 strains of pathogenic viruses were detected within the tick sample from autumn.
Not all ticks in an infective area are infected with transmittable pathogens, and both attachment of the tick and a long feeding session are necessary for diseases to be transmitted. Consequently, tick bites often do not lead to infection, especially if the ticks are removed within 36 hours. Adult ticks can be removed with fine-tipped tweezers or proprietary tick removal tools, before then disinfecting the wound. In Australia and New Zealand, where tick-borne infections are less common than tick reactions, the Australasian Society of Clinical Immunology and Allergy recommends seeking medical assistance or killing ticks in-situ by freezing and then leaving them to fall out to prevent allergic/anaphylactic reactions. Professor Sheryl van Nunen, whose research in 2007 identified tick-induced mammalian meat allergy, famously said "tweezers are tick squeezers", referring to the tick toxins squeezed into people attempting to remove ticks with tweezers. Ticks can be disposed of by flushing them down the toilet, placing them in a container of soapy water or alcohol, or sticking them to tape that can then be folded over and thrown away.
Bifenthrin and permethrin, both pyrethroids, are sometimes used as tick-control measures, although they have the disadvantage of being carcinogenic and able to attack the nervous systems of other species besides ticks. Those who walk through tick-infested areas can make it harder for ticks to latch onto them by tucking their trousers into boots made of smooth rubber, which ticks have trouble climbing.
Research since 2008 has documented red-meat allergies (mammalian meat allergy and Alpha-gal allergy) in the U.S. due to lone star tick bites. The range of the problem has been expanding with the range of the tick. Other species of ticks are known for being responsible for meat allergies in other countries, including Sweden, Germany, and Australia.
Many tick-transmitted viruses, such as Crimean–Congo hemorrhagic fever virus, Kyasanur Forest disease virus, Alkhumra hemorrhagic fever virus, and Omsk hemorrhagic fever virus, are classified as dangerous enough to require biosafety level 4 precautions in laboratory environments. This includes five levels of containment, viz., storage vials within humidified desiccators, within environmental chambers, within a tick suite, within a BSL4 laboratory. Precautions such as glove boxes, sticky pads, Vaseline barriers, safety suits, gloves, sticky tape, silicone vacuum grease, sticky trap paste, and micro mesh are used to safely contain ticks and prevent them from escaping.
Population control measures
With the possible exception of widespread DDT use in the Soviet Union, attempts to limit the population or distribution of disease-causing ticks have been quite unsuccessful.
The parasitoid encyrtid wasp Ixodiphagus hookeri has been investigated for its potential to control tick populations. It lays its eggs into ticks; the hatching wasps kill their hosts.
Predators and competitors of tick hosts can indirectly reduce the density of infected nymphs, thereby lowering tick-borne disease risk by lowering the density and/or tick burden of reservoir-competent hosts. A study in the Netherlands found that the number of larval ticks on bank voles and wood mice was lower at sites with significant red fox (Vulpes vulpes) and stone marten (Martes foina) activity.
This supports the results of a study from the northeastern United States, in which the incidence of Lyme borreliosis was negatively correlated with the density of red fox, possibly because foxes decrease the density of white-footed mice (Peromyscus leucopus), the most important reservoir-competent host for Borrelia burgdorferi.
Another natural form of control for ticks is the helmeted guineafowl, a bird species that consumes mass quantities of ticks. Opossums groom themselves, swallowing many ticks; they are net destroyers of ticks, killing around ninety percent of the ticks that attempt to feed on them. More generally, high animal diversity has a strongly protective effect against tick-borne disease.
Topical tick medicines may be toxic to animals and humans. The synthetic pyrethroid insecticide phenothrin in combination with the hormone analogue methoprene was a popular topical flea and tick therapy for felines. Phenothrin kills adult ticks, while methoprene kills eggs. Some products were withdrawn, and others are known to cause adverse reactions.
| Biology and health sciences | Arachnids | null |
172274 | https://en.wikipedia.org/wiki/Centaur%20%28small%20Solar%20System%20body%29 | Centaur (small Solar System body) | In planetary astronomy, a centaur is a small Solar System body that orbits the Sun between Jupiter and Neptune and crosses the orbits of one or more of the giant planets. Centaurs generally have unstable orbits because of this; almost all their orbits have dynamic lifetimes of only a few million years, but there is one
known centaur, 514107 Kaʻepaokaʻawela, which may be in a stable (though retrograde) orbit. Centaurs typically exhibit the characteristics of both asteroids and comets. They are named after the mythological centaurs that were a mixture of horse and human. Observational bias toward large objects makes determination of the total centaur population difficult. Estimates for the number of centaurs in the Solar System more than 1 km in diameter range from as low as 44,000 to more than 10,000,000.
The first centaur to be discovered, under the definition of the Jet Propulsion Laboratory and the one used here, was 944 Hidalgo in 1920. However, they were not recognized as a distinct population until the discovery of 2060 Chiron in 1977. The largest confirmed centaur is 10199 Chariklo, which at 260 kilometers in diameter is as big as a mid-sized main-belt asteroid, and is known to have a system of rings. It was discovered in 1997.
No centaur has been photographed up close, although there is evidence that Saturn's moon Phoebe, imaged by the Cassini probe in 2004, may be a captured centaur that originated in the Kuiper belt. In addition, the Hubble Space Telescope has gleaned some information about the surface features of 8405 Asbolus.
Ceres may have originated in the region of the outer planets, and if so might be considered an ex-centaur, but the centaurs seen today all originated elsewhere.
Of the objects known to occupy centaur-like orbits, approximately 30 have been found to display comet-like dust comas, with three, 2060 Chiron, 60558 Echeclus, and 29P/Schwassmann-Wachmann 1, having detectable levels of volatile production in orbits entirely beyond Jupiter. Chiron and Echeclus are therefore classified as both centaurs and comets, while Schwassmann-Wachmann 1 has always held a comet designation. Other centaurs, such as 52872 Okyrhoe, are suspected of having shown comas. Any centaur that is perturbed close enough to the Sun is expected to become a comet.
Classification
A centaur has either a perihelion or a semi-major axis between those of the outer planets (between Jupiter and Neptune). Due to the inherent long-term instability of orbits in this region, even centaurs such as and , which do not currently cross the orbit of any planet, are in gradually changing orbits that will be perturbed until they start to cross the orbit of one or more of the giant planets. Some astronomers count only bodies with semimajor axes in the region of the outer planets to be centaurs; others accept any body with a perihelion in the region, as their orbits are similarly unstable.
Discrepant criteria
However, different institutions have different criteria for classifying borderline objects, based on particular values of their orbital elements:
The Minor Planet Center (MPC) defines centaurs as having a perihelion beyond the orbit of Jupiter () and a semi-major axis less than that of Neptune (). Though nowadays the MPC often lists centaurs and scattered disc objects together as a single group.
The Jet Propulsion Laboratory (JPL) similarly defines centaurs as having a semi-major axis, a, between those of Jupiter and Neptune ().
In contrast, the Deep Ecliptic Survey (DES) defines centaurs using a dynamical classification scheme. These classifications are based on the simulated change in behavior of the present orbit when extended over 10 million years. The DES defines centaurs as non-resonant objects whose instantaneous (osculating) perihelia are less than the osculating semi-major axis of Neptune at any time during the simulation. This definition is intended to be synonymous with planet-crossing orbits and to suggest comparatively short lifetimes in the current orbit.
The collection The Solar System Beyond Neptune (2008) defines objects with a semi-major axis between those of Jupiter and Neptune and a Jupiter-relative Tisserand's parameter above 3.05 as centaurs, classifying the objects with a Jupiter-relative Tisserand's parameter below this and, to exclude Kuiper belt objects, an arbitrary perihelion cut-off half-way to Saturn () as Jupiter-family comets, and classifying those objects on unstable orbits with a semi-major axis larger than Neptune's as members of the scattered disc.
Other astronomers prefer to define centaurs as objects that are non-resonant with a perihelion inside the orbit of Neptune that can be shown to likely cross the Hill sphere of a gas giant within the next 10 million years, so that centaurs can be thought of as objects scattered inwards and that interact more strongly and scatter more quickly than typical scattered-disc objects.
The JPL Small-Body Database lists 452 centaurs. There are an additional 116 trans-Neptunian objects (objects with a semi-major axis further than Neptune's, i.e. ) with a perihelion closer than the orbit of Uranus ().
Ambiguous objects
The Gladman & Marsden (2008) criteria would make some objects Jupiter-family comets: Both Echeclus (, ) and Okyrhoe (; ) have traditionally been classified as centaurs. Traditionally considered an asteroid, but classified as a centaur by JPL, Hidalgo (; ) would also change category to a Jupiter-family comet. Schwassmann-Wachmann 1 (; ) has been categorized as both a centaur and a Jupiter-family comet depending on the definition used.
Other objects caught between these differences in classification methods include , which has a semi-major axis of 32 AU but crosses the orbits of both Uranus and Neptune. It is listed as an outer centaur by the Deep Ecliptic Survey (DES). Among the inner centaurs, (434620) 2005 VD, with a perihelion distance very near Jupiter, is listed as a centaur by both JPL and DES.
A recent orbital simulation of the evolution of Kuiper Belt Objects through the centaur region has identified a short-lived "orbital gateway" between 5.4 and 7.8 AU through which 21% of all centaurs pass, including 72% of the centaurs that become Jupiter-family comets. Four objects are known to occupy this region, including 29P/Schwassmann-Wachmann, P/2010 TO20 LINEAR-Grauer, P/2008 CL94 Lemmon, and 2016 LN8, but the simulations indicate that there may of order 1000 more objects >1 km in radius that have yet to be detected. Objects in this gateway region can display significant activity and are in an important evolutionary transition state that further blurs the distinction between the centaur and Jupiter-family comet populations.
The Committee on Small Body Nomenclature of the International Astronomical Union has not formally weighed in on any side of the debate. Instead, it has adopted the following naming convention for such objects: Befitting their centaur-like transitional orbits between TNOs and comets, "objects on unstable, non-resonant, giant-planet-crossing orbits with semimajor axes greater than Neptune's" are to be named for other hybrid and shape-shifting mythical creatures. Thus far, only the binary objects Ceto and Phorcys and Typhon and Echidna have been named according to the new policy.
Centaurs with measured diameters listed as possible dwarf planets according to Mike Brown's website include 10199 Chariklo, and 2060 Chiron.
Orbits
Distribution
The diagram illustrates the orbits of known centaurs in relation to the orbits of the planets. For selected objects, the eccentricity of the orbits is represented by red segments (extending from perihelion to aphelion).
The orbits of centaurs show a wide range of eccentricity, from highly eccentric (Pholus, Asbolus, Amycus, Nessus) to more circular (Chariklo and the Saturn-crossers Thereus and Okyrhoe).
To illustrate the range of the orbits' parameters, the diagram shows a few objects with very unusual orbits, plotted in yellow :
(Apollo asteroid) follows an extremely eccentric orbit (), leading it from inside Earth's orbit (0.94 AU) to well beyond Neptune ()
follows a quasi-circular orbit ()
has the lowest inclination ().
is one of a small proportion of centaurs with an extreme prograde inclination (). It follows such a highly inclined orbit (79°) that, while it crosses from the distance of the asteroid belt from the Sun to past the distance of Saturn, if its orbit is projected onto the plane of Jupiter's orbit, it does not even go out as far as Jupiter.
Over a dozen known centaurs follow retrograde orbits. Their inclinations range from modest (e.g., 160° for Dioretsa) to extreme (; e.g. 105° for ).
Seventeen of these high-inclination, retrograde centaurs were controversially claimed to have an interstellar origin.
Changing orbits
Because the centaurs are not protected by orbital resonances, their orbits are unstable within a timescale of 106–107 years. For example, 55576 Amycus is in an unstable orbit near the 3:4 resonance of Uranus. Dynamical studies of their orbits indicate that being a centaur is probably an intermediate orbital state of objects transitioning from the Kuiper belt to the Jupiter family of short-period comets. (679997) 2023 RB will have its orbit notably changed by a close approach to Saturn in 2201.
Objects may be perturbed from the Kuiper belt, whereupon they become Neptune-crossing and interact gravitationally with that planet (see theories of origin). They then become classed as centaurs, but their orbits are chaotic, evolving relatively rapidly as the centaur makes repeated close approaches to one or more of the outer planets. Some centaurs will evolve into Jupiter-crossing orbits whereupon their perihelia may become reduced into the inner Solar System and they may be reclassified as active comets in the Jupiter family if they display cometary activity. Centaurs will thus ultimately collide with the Sun or a planet or else they may be ejected into interstellar space after a close approach to one of the planets, particularly Jupiter.
Physical characteristics
Compared to dwarf planets and asteroids, the relatively small size and distance of centaurs precludes remote observation of surfaces, but colour indices and spectra can provide clues about surface composition and insight into the origin of the bodies.
Colours
The colours of centaurs are very diverse, which challenges any simple model of surface composition. In the side-diagram, the colour indices are measures of apparent magnitude of an object through blue (B), visible (V) (i.e. green-yellow) and red (R) filters. The diagram illustrates these differences (in exaggerated colours) for all centaurs with known colour indices. For reference, two moons: Triton and Phoebe, and planet Mars are plotted (yellow labels, size not to scale).
Centaurs appear to be grouped into two classes:
very red – for example 5145 Pholus
blue (or blue-grey, according to some authors) – for example 2060 Chiron or
There are numerous theories to explain this colour difference, but they can be broadly divided into two categories:
The colour difference results from a difference in the origin and/or composition of the centaur (see origin below)
The colour difference reflects a different level of space-weathering from radiation and/or cometary activity.
As examples of the second category, the reddish colour of Pholus has been explained as a possible mantle of irradiated red organics, whereas Chiron has instead had its ice exposed due to its periodic cometary activity, giving it a blue/grey index. The correlation with activity and color is not certain, however, as the active centaurs span the range of colors from blue (Chiron) to red (166P/NEAT). Alternatively, Pholus may have been only recently expelled from the Kuiper belt, so that surface transformation processes have not yet taken place.
Delsanti et al. suggest multiple competing processes: reddening by the radiation, and blushing by collisions.
Spectra
The interpretation of spectra is often ambiguous, related to particle sizes and other factors, but the spectra offer an insight into surface composition. As with the colours, the observed spectra can fit a number of models of the surface.
Water ice signatures have been confirmed on a number of centaurs (including 2060 Chiron, 10199 Chariklo and 5145 Pholus). In addition to the water ice signature, a number of other models have been put forward:
Chariklo's surface has been suggested to be a mixture of tholins (like those detected on Titan and Triton) with amorphous carbon.
Pholus has been suggested to be covered by a mixture of Titan-like tholins, carbon black, olivine and methanol ice.
The surface of 52872 Okyrhoe has been suggested to be a mixture of kerogens, olivines and a small percentage of water ice.
8405 Asbolus has been suggested to be a mixture of 15% Triton-like tholins, 8% Titan-like tholin, 37% amorphous carbon and 40% ice tholin.
Chiron appears to be the most complex. The spectra observed vary depending on the period of the observation. Water ice signature was detected during a period of low activity and disappeared during high activity.
Similarities to comets
Observations of Chiron in 1988 and 1989 near its perihelion found it to display a coma (a cloud of gas and dust evaporating from its surface). It is thus now officially classified as both a minor planet and a comet, although it is far larger than a typical comet and there is some lingering controversy. Other centaurs are being monitored for comet-like activity: so far two, 60558 Echeclus, and 166P/NEAT have shown such behavior. 166P/NEAT was discovered while it exhibited a coma, and so is classified as a comet, though its orbit is that of a centaur. 60558 Echeclus was discovered without a coma but recently became active, and so it too is now classified as both a comet and an asteroid. Overall, there are ~30 centaurs for which activity has been detected, with the active population biased toward objects with smaller perihelion distances.
Carbon monoxide has been detected in 60558 Echeclus
and Chiron
in very small amounts, and the derived CO production rate was calculated to be sufficient to account for the observed coma. The calculated CO production rate from both 60558 Echeclus and Chiron is substantially lower than what is typically observed for 29P/Schwassmann–Wachmann, another distantly active comet often classified as a centaur.
There is no clear orbital distinction between centaurs and comets. Both 29P/Schwassmann-Wachmann and 39P/Oterma have been referred to as centaurs since they have typical centaur orbits. The comet 39P/Oterma is currently inactive and was seen to be active only before it was perturbed into a centaur orbit by Jupiter in 1963. The faint comet 38P/Stephan–Oterma would probably not show a coma if it had a perihelion distance beyond Jupiter's orbit at 5 AU. By the year 2200, comet 78P/Gehrels will probably migrate outwards into a centaur-like orbit.
Rotational periods
A periodogram analysis of the light-curves of these Chiron and Chariklo gives respectively the following rotational periods: 5.5±0.4~h and 7.0± 0.6~h.
Size, density, reflectivity
Centaurs can reach diameters up to hundreds of kilometers. The largest centaurs have diameters in excess of 300 km, and primarily reside beyond 20 AU.
Hypotheses of origin
The study of centaurs’ origins is rich in recent developments, but any conclusions are still hampered by limited physical data. Different models have been put forward for possible origin of centaurs.
Simulations indicate that the orbit of some Kuiper belt objects can be perturbed, resulting in the object's expulsion so that it becomes a centaur. Scattered disc objects would be dynamically the best candidates (For instance, the centaurs could be part of an "inner" scattered disc of objects perturbed inwards from the Kuiper belt.) for such expulsions, but their colours do not fit the bicoloured nature of the centaurs. Plutinos are a class of Kuiper belt object that display a similar bicoloured nature, and there are suggestions that not all plutinos' orbits are as stable as initially thought, due to perturbation by Pluto.
Further developments are expected with more physical data on Kuiper belt objects.
Some centaurs may have their origin in fragmentation episodes, perhaps triggered during close encounters with Jupiter. The orbits of centaurs 2020 MK4, P/2008 CL94 (Lemmon), and P/2010 TO20 (LINEAR-Grauer) pass close to that of comet 29P/Schwassmann–Wachmann, the first discovered centaur and close encounters are possible in which one of the objects traverses the coma of 29P when active.
At least one centaur, 2013 VZ70, might have an origin among Saturn's irregular moon population via impact, fragmentation, or tidal disruption.
Notable centaurs
| Physical sciences | Planetary science | Astronomy |
172321 | https://en.wikipedia.org/wiki/Hand%2C%20foot%2C%20and%20mouth%20disease | Hand, foot, and mouth disease | Hand, foot, and mouth disease (HFMD) is a common infection caused by a group of enteroviruses. It typically begins with a fever and feeling generally unwell. This is followed a day or two later by flat discolored spots or bumps that may blister, on the hands, feet and mouth and occasionally buttocks and groin. Signs and symptoms normally appear 3–6 days after exposure to the virus. The rash generally resolves on its own in about a week.
The viruses that cause HFMD are spread through close personal contact, through the air from coughing, and via the feces of an infected person. Contaminated objects can also spread the disease. Coxsackievirus A16 is the most common cause, and enterovirus 71 is the second-most common cause. Other strains of coxsackievirus and enterovirus can also be responsible. Some people may carry and pass on the virus despite having no symptoms of disease. Other animals are not involved. Diagnosis can often be made based on symptoms. Occasionally, a throat or stool sample may be tested for the virus.
Most people with hand, foot, and mouth disease get better on their own in 7 to 10 days. Most cases require no specific treatment. No antiviral medication or vaccine is available, but development efforts are underway. For fever and for painful mouth sores, over-the-counter pain medications such as ibuprofen may be used, though aspirin should be avoided in children. The illness is usually not serious. Occasionally, intravenous fluids are given to children who are dehydrated. Very rarely, viral meningitis or encephalitis may complicate the disease. Because HFMD is normally mild, some jurisdictions allow children to continue to go to child care and schools as long as they have no fever or uncontrolled drooling with mouth sores, and as long as they feel well enough to participate in classroom activities.
HFMD occurs in all areas of the world. It often occurs in small outbreaks in nursery schools or kindergartens. Large outbreaks have been occurring in Asia since 1997. It usually occurs during the spring, summer, and fall months. Typically it occurs in children less than five years old but can occasionally occur in adults. HFMD should not be confused with foot-and-mouth disease (also known as hoof-and-mouth disease), which mostly affects livestock.
Signs and symptoms
Common constitutional signs and symptoms of HFMD include fever, nausea, vomiting, feeling tired, generalized discomfort, loss of appetite, and irritability in infants and toddlers. Skin lesions frequently develop in the form of a rash of flat discolored spots and bumps which may be followed by vesicular sores with blisters on palms of the hands, soles of the feet, buttocks, and sometimes on the lips. The rash is rarely itchy for children, but can be extremely itchy for adults. Painful facial ulcers, blisters, or lesions may also develop in or around the nose or mouth. HFMD usually resolves on its own after 7–10 days. Most cases of the disease are relatively harmless, but complications including encephalitis, meningitis, and paralysis that mimics the neurological symptoms of polio can occur.
Cause
The viruses that cause the disease are of the Picornaviridae family. Coxsackievirus A16 is the most common cause of HFMD. Enterovirus 71 (EV-71) is the second-most common cause. Many other strains of coxsackievirus and enterovirus can also be responsible.
Transmission
HFMD is highly contagious and is transmitted by nasopharyngeal secretions such as saliva or nasal mucus, by direct contact, or by fecal–oral transmission. It is possible to be infectious for days to weeks after the symptoms have resolved.
Childcare settings are the most common places for HFMD to be contracted because of toilet training, diaper changes, and children's propensity to put their hands into their mouths. HFMD is contracted through nose and throat secretions such as saliva, sputum, and nasal mucus as well as fluid in blisters, and stool.
Diagnosis
A diagnosis usually can be made by the presenting signs and symptoms alone. If the diagnosis is unclear, a throat swab or stool specimen may be taken to identify the virus by culture. The common incubation period (the time between infection and onset of symptoms) ranges from three to six days. Early detection of HFMD is important in preventing an outbreak in the pediatric population.
Prevention
Preventive measures include avoiding direct contact with infected individuals (including keeping infected children home from school), proper cleaning of shared utensils, disinfecting contaminated surfaces, and proper hand hygiene. These measures are effective in decreasing the transmission of the viruses responsible for HFMD.
Protective habits include hand washing and disinfecting surfaces in play areas. Breastfeeding has also been shown to decrease rates of severe HFMD, though does not reduce the risk of the infection of the disease.
Vaccine
A vaccine known as the EV71 vaccine is available to prevent HFMD in China . No vaccine is currently available in the United States.
Treatment
Medications are usually not needed as hand, foot, and mouth disease is a viral disease that typically resolves on its own. Currently, there is no specific curative treatment for hand, foot, and mouth disease. Disease management typically focuses on achieving symptomatic relief. Pain from the sores may be eased with the use of analgesic medications. Infection in older children, adolescents, and adults is typically mild and lasts approximately 1 week, but may occasionally run a longer course. Fever reducers can help decrease body temperature.
A minority of individuals with hand, foot, and mouth disease may require hospital admission due to complications such as inflammation of the brain, inflammation of the meninges, or acute flaccid paralysis. Non-neurologic complications such as inflammation of the heart, fluid in the lungs, or bleeding into the lungs may also occur.
Complications
Complications from the viral infections that cause HFMD are rare but require immediate medical treatment if present. HFMD infections caused by Enterovirus 71 tend to be more severe and are more likely to have neurologic or cardiac complications including death than infections caused by Coxsackievirus A16. Viral or aseptic meningitis can occur with HFMD in rare cases and is characterized by fever, headache, stiff neck, or back pain. The condition is usually mild and clears without treatment; however, hospitalization for a short time may be needed. Other serious complications of HFMD include encephalitis (inflammation of the brain), or flaccid paralysis in rare circumstances.
Fingernail and toenail loss have been reported in children 4–8 weeks after having HFMD. The relationship between HFMD and the reported nail loss is unclear; however, it is temporary and nail growth resumes without treatment.
Minor complications due to symptoms can occur such as dehydration, due to mouth sores causing discomfort with intake of foods and fluid.
Epidemiology
Hand, foot and mouth disease most commonly occurs in children under the age of 10 and more often under the age of 5, but it can also affect adults with varying symptoms. It tends to occur in outbreaks during the spring, summer, and autumn seasons. This is believed to be due to heat and humidity improving spread. HFMD is more common in rural areas than urban areas; however, socioeconomic status and hygiene levels need to be considered. Poor hygiene is a risk factor for HFMD.
Outbreaks
In 1997, an outbreak occurred in Sarawak of Malaysia with 600 cases and over 30 children died.
In 1998, there was an outbreak in Taiwan, affecting mainly children. There were 405 severe complications, and 78 children died. The total number of cases in that epidemic is estimated to have been 1.5 million.
In 2008 an outbreak in China, beginning in March in Fuyang, Anhui, led to 25,000 infections, and 42 deaths, by May 13. Similar outbreaks were reported in Singapore (more than 2,600 cases as of April 20, 2008), Vietnam (2,300 cases, 11 deaths), Mongolia (1,600 cases), and Brunei (1053 cases from June–August 2008)
In 2009 17 children died in an outbreak during March and April 2009 in China's eastern Shandong Province, and 18 children died in the neighboring Henan Province. Out of 115,000 reported cases in China from January to April, 773 were severe and 50 were fatal.
In 2010 in China, an outbreak occurred in southern China's Guangxi Autonomous Region as well as Guangdong, Henan, Hebei, and Shandong provinces. Until March, 70,756 children were infected and 40 died from the disease. By June, the peak season for the disease, 537 had died.
The World Health Organization reporting between January and October 2011 (1,340,259) states the number of cases in China had dropped by approx 300,000 from 2010 (1,654,866) cases, with new cases peaking in June. There were 437 deaths, down from 2010 (537 deaths).
In December 2011, the California Department of Public Health identified a strong form of the virus, coxsackievirus A6 (CVA6), where nail loss in children is common.
In 2012 in Alabama, United States there was an outbreak of an unusual type of the disease. It occurred in a season when it is not usually seen and affected teenagers and older adults. There were some hospitalizations due to the disease but no reported deaths.
In 2012 in Cambodia, 52 of 59 reviewed cases of children reportedly dead () due to a mysterious disease were diagnosed to be caused by a virulent form of HFMD. Although a significant degree of uncertainty exists with reference to the diagnosis, the WHO report states, "Based on the latest laboratory results, a significant proportion of the samples tested positive for enterovirus 71 (EV-71), which causes hand foot and mouth disease (HFMD). The EV-71 virus has been known to generally cause severe complications amongst some patients."
HFMD infected 1,520,274 people with up to 431 deaths reported at the end of July in 2012 in China.
In 2018, more than 50,000 cases occurred through a nationwide outbreak in Malaysia with two deaths also reported.
India 2022
An outbreak of an illness referred to as tomato fever or tomato flu was identified in the Kollam district on May 6, 2022. The illness is endemic to Kerala, India and gets its name because of the red and round blisters it causes, which look like tomatoes. The disease may be a new variant of the viral HFMD or an effect of chikungunya or dengue fever. Flu may be a misnomer.
The condition mainly affects children under the age of five. An article in The Lancet states that the appearance of the blisters is similar to that seen in Mpox, and the illness is not thought to be related to SARS-CoV-2. Symptoms, treatment and prevention are similar to HFMD.
History
HFMD cases were first described clinically in Canada and New Zealand in 1957. The disease was termed "Hand Foot and Mouth Disease", by Thomas Henry Flewett, after a similar outbreak in 1960.
Research
Novel antiviral agents to prevent and treat infection with the viruses responsible for HFMD are currently under development. Preliminary studies have shown inhibitors of the EV-71 viral capsid to have potent antiviral activity.
| Biology and health sciences | Infectious disease | null |
172323 | https://en.wikipedia.org/wiki/Rubella | Rubella | Rubella, also known as German measles or three-day measles, is an infection caused by the rubella virus. This disease is often mild, with half of people not realizing that they are infected. A rash may start around two weeks after exposure and last for three days. It usually starts on the face and spreads to the rest of the body. The rash is sometimes itchy and is not as bright as that of measles. Swollen lymph nodes are common and may last a few weeks. A fever, sore throat, and fatigue may also occur. Joint pain is common in adults. Complications may include bleeding problems, testicular swelling, encephalitis, and inflammation of nerves. Infection during early pregnancy may result in a miscarriage or a child born with congenital rubella syndrome (CRS). Symptoms of CRS manifest as problems with the eyes such as cataracts, deafness, as well as affecting the heart and brain. Problems are rare after the 20th week of pregnancy.
Rubella is usually spread from one person to the next through the air via coughs of people who are infected. People are infectious during the week before and after the appearance of the rash. Babies with CRS may spread the virus for more than a year. Only humans are infected. Insects do not spread the disease. Once recovered, people are immune to future infections. Testing is available that can verify immunity. Diagnosis is confirmed by finding the virus in the blood, throat, or urine. Testing the blood for antibodies may also be useful.
Rubella is preventable with the rubella vaccine, with a single dose being more than 95% effective. Often it is given in combination with the measles vaccine and mumps vaccine, known as the MMR vaccine. When some, but less than 80%, of a population is vaccinated, more women may reach childbearing age without developing immunity by infection or vaccination, thus possibly raising CRS rates. Once infected there is no specific treatment.
Rubella is a common infection in many areas of the world. Each year about 100,000 cases of congenital rubella syndrome occur. Rates of disease have decreased in many areas as a result of vaccination. There are ongoing efforts to eliminate the disease globally. In April 2015, the World Health Organization declared the Americas free of rubella transmission. The name "rubella" is from Latin and means little red. It was first described as a separate disease by German physicians in 1814, resulting in the name "German measles".
Signs and symptoms
Rubella has symptoms similar to those of flu. However, the primary symptom of rubella virus infection is the appearance of a rash (exanthem) on the face which spreads to the trunk and limbs and usually fades after three days, which is why it is often referred to as three-day measles. The facial rash usually clears as it spreads to other parts of the body. Other symptoms include low-grade fever, swollen glands (sub-occipital and posterior cervical lymphadenopathy), joint pains, headache, and conjunctivitis.
The swollen glands or lymph nodes can persist for up to a week and the fever rarely rises above 38 °C (100.4 °F). The rash of rubella is typically pink or light red. The rash causes itching and often lasts for about three days. The rash disappears after a few days with no staining or peeling of the skin. When the rash clears up, the skin might shed in very small flakes where the rash covered it. Forchheimer spots occur in 20% of cases and are characterized by small, red papules on the area of the soft palate.
Rubella can affect anyone of any age. Adult females are particularly prone to arthritis and joint pains.
In children, rubella normally causes symptoms that last two days and include:
Rash begins on the face which spreads to the rest of the body.
Low fever of less than .
Posterior cervical lymphadenopathy.
In older children and adults, additional symptoms may be present, including
Swollen glands
Coryza (cold-like symptoms)
Aching joints (especially in young females)
Severe complications of rubella include:
Brain inflammation (encephalitis)
Low platelet count
Ear infection
Coryza in rubella may convert to pneumonia, either direct viral pneumonia or secondary bacterial pneumonia, and bronchitis (either viral bronchitis or secondary bacterial bronchitis).
Congenital rubella syndrome
Rubella can cause congenital rubella syndrome in the newborn, this being the most severe sequela of rubella. The syndrome (CRS) follows intrauterine infection by the rubella virus and comprises cardiac, cerebral, ophthalmic, and auditory defects. It may also cause prematurity, low birth weight, neonatal thrombocytopenia, anemia, and hepatitis. The risk of major defects in organogenesis is highest for infection in the first trimester. CRS is the main reason a vaccine for rubella was developed.
80–90% of mothers who contract rubella within the critical first trimester have either a miscarriage or a stillborn baby. If the fetus survives the infection, it can be born with severe heart disorders (patent ductus arteriosus being the most common), blindness, deafness, or other life-threatening organ disorders. The skin manifestations are called "blueberry muffin lesions". For these reasons, rubella is included in the TORCH complex of perinatal infections.
About 100,000 cases of this condition occur each year.
Cause
The disease is caused by the rubella virus, in the genus Rubivirus from the family Matonaviridae, that is enveloped and has a single-stranded RNA genome. The virus is transmitted by the respiratory route and replicates in the nasopharynx and lymph nodes. The virus is found in the blood 5 to 7 days after infection and spreads throughout the body. The virus has teratogenic properties and is capable of crossing the placenta and infecting the fetus where it stops cells from developing or destroys them. During this incubation period, the patient is contagious typically for about one week before he/she develops a rash and for about one week thereafter.
Increased susceptibility to infection might be inherited as there is some indication that HLA-A1 or factors surrounding A1 on extended haplotypes are involved in virus infection or non-resolution of the disease.
Diagnosis
Rubella virus specific IgM antibodies are present in people recently infected by rubella virus, but these antibodies can persist for over a year, and a positive test result needs to be interpreted with caution. The presence of these antibodies along with, or a short time after, the characteristic rash confirms the diagnosis.
Prevention
Rubella infections are prevented by active immunization programs using live attenuated virus vaccines. Two live attenuated virus vaccines, RA 27/3 and Cendehill strains, were effective in the prevention of adult disease. However, their use in prepubertal females did not produce a significant fall in the overall incidence rate of CRS in the UK. Reductions were only achieved by immunisation of all children.
The vaccine is now usually given as part of the MMR vaccine. The WHO recommends the first dose be given at 12 to 18 months of age with a second dose at 36 months. Pregnant women are usually tested for immunity to rubella early on. Women found to be susceptible are not vaccinated until after the baby is born because the vaccine contains live virus.
The immunisation program has been quite successful. Cuba declared the disease eliminated in the 1990s, and in 2004 the Centers for Disease Control and Prevention announced that both the congenital and acquired forms of rubella had been eliminated from the United States. The World Health Organization declared Australia rubella free in October 2018.
Screening for rubella susceptibility by history of vaccination or by serology is recommended in the United States for all women of childbearing age at their first preconception counseling visit to reduce incidence of congenital rubella syndrome (CRS). It is recommended that all susceptible non-pregnant women of childbearing age should be offered rubella vaccination. Due to concerns about possible teratogenicity, use of MMR vaccine is not recommended during pregnancy. Instead, susceptible pregnant women should be vaccinated as soon as possible in the postpartum period.
In susceptible people passive immunization, in the form of polyclonal immunoglobulins, appears effective up to the fifth day post-exposure.
Treatment
There is no specific treatment for rubella; however, management is a matter of responding to symptoms to diminish discomfort. Treatment of newborn babies is focused on management of the complications. Congenital heart defects and cataracts can be corrected by direct surgery.
Management for ocular congenital rubella syndrome (CRS) is similar to that for age-related macular degeneration, including counseling, regular monitoring, and the provision of low vision devices, if required.
Prognosis
Rubella infection of children and adults is usually mild, self-limiting, and often asymptomatic. The prognosis in children born with CRS is poor.
Epidemiology
Rubella occurs worldwide. The virus tends to peak during the spring in countries with temperate climates. Before the vaccine against rubella was introduced in 1969, widespread outbreaks usually occurred every 6–9 years in the United States and 3–5 years in Europe, mostly affecting children in the 5–9 year old age group. Since the introduction of vaccine, occurrences have become rare in those countries with high uptake rates.
Vaccination has interrupted the transmission of rubella in the Americas: no endemic case has been observed since February 2009. Vaccination is still strongly recommended as the virus could be reintroduced from other continents should vaccination rates in the Americas drop. During the epidemic in the US between 1962 and 1965, rubella virus infections during pregnancy were estimated to have caused 30,000 stillbirths and 20,000 children to be born impaired or disabled as a result of CRS. Universal immunisation producing a high level of herd immunity is important in the control of epidemics of rubella.
In the UK, there remains a large population of men susceptible to rubella who have not been vaccinated. Outbreaks of rubella occurred amongst many young men in the UK in 1993 and in 1996 the infection was transmitted to pregnant women, many of whom were immigrants and were susceptible. Outbreaks still arise, usually in developing countries where the vaccine is not as accessible. The complications encountered in pregnancy from rubella infection (miscarriage, fetal death, congenital rubella syndrome) are more common in Africa and Southeast Asia at a rate of 121 per 100,000 live births compared to 2 per 100,000 live births in the Americas and Europe.
In Japan, 15,000 cases of rubella and 43 cases of congenital rubella syndrome were reported to the National Epidemiological Surveillance of Infectious Diseases between October 15, 2012, and March 2, 2014, during the 2012–13 rubella outbreak in Japan. They mainly occurred in men aged 31–51 and young adults aged 24–34.
History
Rubella was first described in the mid-eighteenth century. German physician and chemist, Friedrich Hoffmann, made the first clinical description of rubella in 1740, which was confirmed by de Bergen in 1752 and Orlow in 1758.
In 1814, George de Maton first suggested that it be considered a disease distinct from both measles and scarlet fever. All these physicians were German, and the disease was known as Rötheln (contemporary German Röteln). (Rötlich means "reddish" or "pink" in German.) The fact that three Germans described it led to the common name of "German measles." Henry Veale, an English Royal Artillery surgeon, described an outbreak in India. He coined the name "rubella" (from the Latin word, meaning "little red") in 1866.
It was formally recognised as an individual entity in 1881, at the International Congress of Medicine in London. In 1914, Alfred Fabian Hess theorised that rubella was caused by a virus, based on work with monkeys. In 1938, Hiro and Tosaka confirmed this by passing the disease to children using filtered nasal washings from acute cases.
In 1940, there was a widespread epidemic of rubella in Australia. Subsequently, ophthalmologist Norman McAllister Gregg found 78 cases of congenital cataracts in infants and 68 of them were born to mothers who had caught rubella in early pregnancy. Gregg published an account, Congenital Cataract Following German Measles in the Mother, in 1941. He described a variety of problems now known as congenital rubella syndrome (CRS) and noticed that the earlier the mother was infected, the worse the damage was. Since no vaccine was yet available, some popular magazines promoted the idea of "German measles parties" for infected children to spread the disease to other children (especially girls) to immunize them for life and protect them from later catching the disease when pregnant. The virus was isolated in tissue culture in 1962 by two separate groups led by physicians Paul Douglas Parkman and Thomas Huckle Weller.
There was a pandemic of rubella between 1962 and 1965, starting in Europe and spreading to the United States. In the years 1964–65, the United States had an estimated 12.5 million rubella cases (1964–1965 rubella epidemic). This led to 11,000 miscarriages or therapeutic abortions and 20,000 cases of congenital rubella syndrome. Of these, 2,100 died as neonates, 12,000 were deaf, 3,580 were blind, and 1,800 were intellectually disabled. In New York alone, CRS affected 1% of all births.
In 1967, the molecular structure of rubella was observed under electron microscopy using antigen-antibody complexes by Jennifer M. Best, June Almeida, J E Banatvala and A P Waterson.
In 1969, a live attenuated virus vaccine was licensed. In the early 1970s, a triple vaccine containing attenuated measles, mumps and rubella (MMR) viruses was introduced. By 2006, confirmed cases in the Americas had dropped below 3000 a year. However, a 2007 outbreak in Argentina, Brazil, and Chile pushed the cases to 13,000 that year.
Eradication efforts
On January 22, 2014, the World Health Organization (WHO) and the Pan American Health Organization declared and certified Colombia free of rubella and became the first Latin American country to eliminate the disease within its borders. On April 29, 2015, the Americas became the first WHO region to officially eradicate the disease. The last non-imported cases occurred in 2009 in Argentina and Brazil. The Pan American Health Organization director remarked, "The fight against rubella has taken more than 15 years, but it has paid off with what I believe will be one of the most important pan-American public health achievements of the 21st Century." The declaration was made after 165 million health records and genetically confirming that all recent cases were caused by known imported strains of the virus. Rubella is still common in some regions of the world and Susan E. Reef, team lead for rubella at the CDC's global immunization division, who joined in the announcement, said there was no chance it would be eradicated worldwide before 2020. Rubella is the third disease to be eliminated from the Western Hemisphere with vaccination after smallpox and polio.
Etymology
From "rubrum" the Latin for "red", rubella means "reddish and small". "German" measles derives from "germanus" which means "similar" in this context.
The name rubella is sometimes confused with rubeola, an alternative name for measles in English-speaking countries; the diseases are unrelated. In some other European languages, like Spanish, rubella and rubeola are synonyms, and rubeola is not an alternative name for measles. Thus, in Spanish, rubeola refers to rubella and sarampión refers to measles.
| Biology and health sciences | Viral diseases | Health |
172333 | https://en.wikipedia.org/wiki/Dispersion%20%28optics%29 | Dispersion (optics) | Dispersion is the phenomenon in which the phase velocity of a wave depends on its frequency. Sometimes the term chromatic dispersion is used to refer to optics specifically, as opposed to wave propagation in general. A medium having this common property may be termed a dispersive medium.
Although the term is used in the field of optics to describe light and other electromagnetic waves, dispersion in the same sense can apply to any sort of wave motion such as acoustic dispersion in the case of sound and seismic waves, and in gravity waves (ocean waves). Within optics, dispersion is a property of telecommunication signals along transmission lines (such as microwaves in coaxial cable) or the pulses of light in optical fiber.
In optics, one important and familiar consequence of dispersion is the change in the angle of refraction of different colors of light, as seen in the spectrum produced by a dispersive prism and in chromatic aberration of lenses. Design of compound achromatic lenses, in which chromatic aberration is largely cancelled, uses a quantification of a glass's dispersion given by its Abbe number V, where lower Abbe numbers correspond to greater dispersion over the visible spectrum. In some applications such as telecommunications, the absolute phase of a wave is often not important but only the propagation of wave packets or "pulses"; in that case one is interested only in variations of group velocity with frequency, so-called group-velocity dispersion.
All common transmission media also vary in attenuation (normalized to transmission length) as a function of frequency, leading to attenuation distortion; this is not dispersion, although sometimes reflections at closely spaced impedance boundaries (e.g. crimped segments in a cable) can produce signal distortion which further aggravates inconsistent transit time as observed across signal bandwidth.
Examples
The most familiar example of dispersion is probably a rainbow, in which dispersion causes the spatial separation of a white light into components of different wavelengths (different colors). However, dispersion also has an effect in many other circumstances: for example, group-velocity dispersion causes pulses to spread in optical fibers, degrading signals over long distances; also, a cancellation between group-velocity dispersion and nonlinear effects leads to soliton waves.
Material and waveguide dispersion
Most often, chromatic dispersion refers to bulk material dispersion, that is, the change in refractive index with optical frequency. However, in a waveguide there is also the phenomenon of waveguide dispersion, in which case a wave's phase velocity in a structure depends on its frequency simply due to the structure's geometry. More generally, "waveguide" dispersion can occur for waves propagating through any inhomogeneous structure (e.g., a photonic crystal), whether or not the waves are confined to some region. In a waveguide, both types of dispersion will generally be present, although they are not strictly additive. For example, in fiber optics the material and waveguide dispersion can effectively cancel each other out to produce a zero-dispersion wavelength, important for fast fiber-optic communication.
Material dispersion in optics
Material dispersion can be a desirable or undesirable effect in optical applications. The dispersion of light by glass prisms is used to construct spectrometers and spectroradiometers. However, in lenses, dispersion causes chromatic aberration, an undesired effect that may degrade images in microscopes, telescopes, and photographic objectives.
The phase velocity v of a wave in a given uniform medium is given by
where c is the speed of light in vacuum, and n is the refractive index of the medium.
In general, the refractive index is some function of the frequency f of the light, thus n = n(f), or alternatively, with respect to the wave's wavelength n = n(λ). The wavelength dependence of a material's refractive index is usually quantified by its Abbe number or its coefficients in an empirical formula such as the Cauchy or Sellmeier equations.
Because of the Kramers–Kronig relations, the wavelength dependence of the real part of the refractive index is related to the material absorption, described by the imaginary part of the refractive index (also called the extinction coefficient). In particular, for non-magnetic materials (μ = μ0), the susceptibility χ that appears in the Kramers–Kronig relations is the electric susceptibility χe = n2 − 1.
The most commonly seen consequence of dispersion in optics is the separation of white light into a color spectrum by a prism. From Snell's law it can be seen that the angle of refraction of light in a prism depends on the refractive index of the prism material. Since that refractive index varies with wavelength, it follows that the angle that the light is refracted by will also vary with wavelength, causing an angular separation of the colors known as angular dispersion.
For visible light, refraction indices n of most transparent materials (e.g., air, glasses) decrease with increasing wavelength λ:
or generally,
In this case, the medium is said to have normal dispersion. Whereas if the index increases with increasing wavelength (which is typically the case in the ultraviolet), the medium is said to have anomalous dispersion.
At the interface of such a material with air or vacuum (index of ~1), Snell's law predicts that light incident at an angle θ to the normal will be refracted at an angle arcsin(). Thus, blue light, with a higher refractive index, will be bent more strongly than red light, resulting in the well-known rainbow pattern.
Group-velocity dispersion
Beyond simply describing a change in the phase velocity over wavelength, a more serious consequence of dispersion in many applications is termed group-velocity dispersion (GVD). While phase velocity v is defined as v = c/n, this describes only one frequency component. When different frequency components are combined, as when considering a signal or a pulse, one is often more interested in the group velocity, which describes the speed at which a pulse or information superimposed on a wave (modulation) propagates. In the accompanying animation, it can be seen that the wave itself (orange-brown) travels at a phase velocity much faster than the speed of the envelope (black), which corresponds to the group velocity. This pulse might be a communications signal, for instance, and its information only travels at the group velocity rate, even though it consists of wavefronts advancing at a faster rate (the phase velocity).
It is possible to calculate the group velocity from the refractive-index curve n(ω) or more directly from the wavenumber k = ωn/c, where ω is the radian frequency ω = 2πf. Whereas one expression for the phase velocity is vp = ω/k, the group velocity can be expressed using the derivative: vg = dω/dk. Or in terms of the phase velocity vp,
When dispersion is present, not only the group velocity is not equal to the phase velocity, but generally it itself varies with wavelength. This is known as group-velocity dispersion and causes a short pulse of light to be broadened, as the different-frequency components within the pulse travel at different velocities. Group-velocity dispersion is quantified as the derivative of the reciprocal of the group velocity with respect to angular frequency, which results in group-velocity dispersion = d2k/dω2.
If a light pulse is propagated through a material with positive group-velocity dispersion, then the shorter-wavelength components travel slower than the longer-wavelength components. The pulse therefore becomes positively chirped, or up-chirped, increasing in frequency with time. On the other hand, if a pulse travels through a material with negative group-velocity dispersion, shorter-wavelength components travel faster than the longer ones, and the pulse becomes negatively chirped, or down-chirped, decreasing in frequency with time.
An everyday example of a negatively chirped signal in the acoustic domain is that of an approaching train hitting deformities on a welded track. The sound caused by the train itself is impulsive and travels much faster in the metal tracks than in air, so that the train can be heard well before it arrives. However, from afar it is not heard as causing impulses, but leads to a distinctive descending chirp, amidst reverberation caused by the complexity of the vibrational modes of the track. Group-velocity dispersion can be heard in that the volume of the sounds stays audible for a surprisingly long time, up to several seconds.
Dispersion control
The result of GVD, whether negative or positive, is ultimately temporal spreading of the pulse. This makes dispersion management extremely important in optical communications systems based on optical fiber, since if dispersion is too high, a group of pulses representing a bit-stream will spread in time and merge, rendering the bit-stream unintelligible. This limits the length of fiber that a signal can be sent down without regeneration. One possible answer to this problem is to send signals down the optical fibre at a wavelength where the GVD is zero (e.g., around 1.3–1.5 μm in silica fibres), so pulses at this wavelength suffer minimal spreading from dispersion. In practice, however, this approach causes more problems than it solves because zero GVD unacceptably amplifies other nonlinear effects (such as four-wave mixing). Another possible option is to use soliton pulses in the regime of negative dispersion, a form of optical pulse which uses a nonlinear optical effect to self-maintain its shape. Solitons have the practical problem, however, that they require a certain power level to be maintained in the pulse for the nonlinear effect to be of the correct strength. Instead, the solution that is currently used in practice is to perform dispersion compensation, typically by matching the fiber with another fiber of opposite-sign dispersion so that the dispersion effects cancel; such compensation is ultimately limited by nonlinear effects such as self-phase modulation, which interact with dispersion to make it very difficult to undo.
Dispersion control is also important in lasers that produce short pulses. The overall dispersion of the optical resonator is a major factor in determining the duration of the pulses emitted by the laser. A pair of prisms can be arranged to produce net negative dispersion, which can be used to balance the usually positive dispersion of the laser medium. Diffraction gratings can also be used to produce dispersive effects; these are often used in high-power laser amplifier systems. Recently, an alternative to prisms and gratings has been developed: chirped mirrors. These dielectric mirrors are coated so that different wavelengths have different penetration lengths, and therefore different group delays. The coating layers can be tailored to achieve a net negative dispersion.
In waveguides
Waveguides are highly dispersive due to their geometry (rather than just to their material composition). Optical fibers are a sort of waveguide for optical frequencies (light) widely used in modern telecommunications systems. The rate at which data can be transported on a single fiber is limited by pulse broadening due to chromatic dispersion among other phenomena.
In general, for a waveguide mode with an angular frequency ω(β) at a propagation constant β (so that the electromagnetic fields in the propagation direction z oscillate proportional to ei(βz−ωt)), the group-velocity dispersion parameter D is defined as
where λ = 2c/ω is the vacuum wavelength, and vg = dω/dβ is the group velocity. This formula generalizes the one in the previous section for homogeneous media and includes both waveguide dispersion and material dispersion. The reason for defining the dispersion in this way is that |D| is the (asymptotic) temporal pulse spreading Δt per unit bandwidth
Δλ per unit distance travelled, commonly reported in ps/(nm⋅km) for optical fibers.
In the case of multi-mode optical fibers, so-called modal dispersion will also lead to pulse broadening. Even in single-mode fibers, pulse broadening can occur as a result of polarization mode dispersion (since there are still two polarization modes). These are not examples of chromatic dispersion, as they are not dependent on the wavelength or bandwidth of the pulses propagated.
Higher-order dispersion over broad bandwidths
When a broad range of frequencies (a broad bandwidth) is present in a single wavepacket, such as in an ultrashort pulse or a chirped pulse or other forms of spread spectrum transmission, it may not be accurate to approximate the dispersion by a constant over the entire bandwidth, and more complex calculations are required to compute effects such as pulse spreading.
In particular, the dispersion parameter D defined above is obtained from only one derivative of the group velocity. Higher derivatives are known as higher-order dispersion. These terms are simply a Taylor series expansion of the dispersion relation β(ω) of the medium or waveguide around some particular frequency. Their effects can be computed via numerical evaluation of Fourier transforms of the waveform, via integration of higher-order slowly varying envelope approximations, by a split-step method (which can use the exact dispersion relation rather than a Taylor series), or by direct simulation of the full Maxwell's equations rather than an approximate envelope equation.
Spatial dispersion
In electromagnetics and optics, the term dispersion generally refers to aforementioned temporal or frequency dispersion. Spatial dispersion refers to the non-local response of the medium to the space; this can be reworded as the wavevector dependence of the permittivity. For an exemplary anisotropic medium, the spatial relation between electric and electric displacement field can be expressed as a convolution:
where the kernel is dielectric response (susceptibility); its indices make it in general a tensor to account for the anisotropy of the medium. Spatial dispersion is negligible in most macroscopic cases, where the scale of variation of is much larger than atomic dimensions, because the dielectric kernel dies out at macroscopic distances. Nevertheless, it can result in non-negligible macroscopic effects, particularly in conducting media such as metals, electrolytes and plasmas. Spatial dispersion also plays role in optical activity and Doppler broadening, as well as in the theory of metamaterials.
In gemology
In the technical terminology of gemology, dispersion is the difference in the refractive index of a material at the B and G (686.7 nm and 430.8 nm) or C and F (656.3 nm and 486.1 nm) Fraunhofer wavelengths, and is meant to express the degree to which a prism cut from the gemstone demonstrates "fire". Fire is a colloquial term used by gemologists to describe a gemstone's dispersive nature or lack thereof. Dispersion is a material property. The amount of fire demonstrated by a given gemstone is a function of the gemstone's facet angles, the polish quality, the lighting environment, the material's refractive index, the saturation of color, and the orientation of the viewer relative to the gemstone.
In imaging
In photographic and microscopic lenses, dispersion causes chromatic aberration, which causes the different colors in the image not to overlap properly. Various techniques have been developed to counteract this, such as the use of achromats, multielement lenses with glasses of different dispersion. They are constructed in such a way that the chromatic aberrations of the different parts cancel out.
Pulsar emissions
Pulsars are spinning neutron stars that emit pulses at very regular intervals ranging from milliseconds to seconds. Astronomers believe that the pulses are emitted simultaneously over a wide range of frequencies. However, as observed on Earth, the components of each pulse emitted at higher radio frequencies arrive before those emitted at lower frequencies. This dispersion occurs because of the ionized component of the interstellar medium, mainly the free electrons, which make the group velocity frequency-dependent. The extra delay added at a frequency is
where the dispersion constant kDM is given by
and the dispersion measure (DM) is the column density of free electrons (total electron content) i.e. the number density of electrons ne integrated along the path traveled by the photon from the pulsar to the Earth and is given by
with units of parsecs per cubic centimetre (1 pc/cm3 = 30.857 m−2).
Typically for astronomical observations, this delay cannot be measured directly, since the emission time is unknown. What can be measured is the difference in arrival times at two different frequencies. The delay Δt between a high-frequency hi and a low-frequency lo component of a pulse will be
Rewriting the above equation in terms of Δt allows one to determine the DM by measuring pulse arrival times at multiple frequencies. This in turn can be used to study the interstellar medium, as well as allow observations of pulsars at different frequencies to be combined.
| Physical sciences | Optics | Physics |
172396 | https://en.wikipedia.org/wiki/Lichen | Lichen | A lichen ( , ) is a hybrid colony of algae or cyanobacteria living symbiotically among filaments of multiple fungi species, along with yeasts and bacteria embedded in the cortex or "skin", in a mutualistic relationship. Lichens are the lifeform that first brought the term symbiosis (as Symbiotismus) under biological context.
Lichens have since been recognized as important actors in nutrient cycling and producers which many higher trophic feeders feed on, such as reindeer, gastropods, nematodes, mites, and springtails. Lichens have properties different from those of their component organisms. They come in many colors, sizes, and forms and are sometimes plant-like, but are not plants. They may have tiny, leafless branches (fruticose); flat leaf-like structures (foliose); grow crust-like, adhering tightly to a surface (substrate) like a thick coat of paint (crustose); have a powder-like appearance (leprose); or other growth forms.
A macrolichen is a lichen that is either bush-like or leafy; all other lichens are termed microlichens. Here, "macro" and "micro" do not refer to size, but to the growth form. Common names for lichens may contain the word moss (e.g., "reindeer moss", "Iceland moss"), and lichens may superficially look like and grow with mosses, but they are not closely related to mosses or any plant. Lichens do not have roots that absorb water and nutrients as plants do, but like plants, they produce their own energy by photosynthesis. Instead, lichen absorb nutrients from rainwater and the air . When they grow on plants, they do not live as parasites, but instead use the plant's surface as a substrate.
Lichens occur from sea level to high alpine elevations, in many environmental conditions, and can grow on almost any surface. They are abundant growing on bark, leaves, mosses, or other lichens and hanging from branches "living on thin air" (epiphytes) in rainforests and in temperate woodland. They grow on rock, walls, gravestones, roofs, exposed soil surfaces, rubber, bones, and in the soil as part of biological soil crusts. Various lichens have adapted to survive in some of the most extreme environments on Earth: arctic tundra, hot dry deserts, rocky coasts, and toxic slag heaps. They can even live inside solid rock, growing between the grains (endolithic).
There are about 20,000 known species. Some lichens have lost the ability to reproduce sexually, yet continue to speciate. They can be seen as being relatively self-contained miniature ecosystems, where the fungi, algae, or cyanobacteria have the potential to engage with other microorganisms in a functioning system that may evolve as an even more complex composite organism. Lichens may be long-lived, with some considered to be among the oldest living things. They are among the first living things to grow on fresh rock exposed after an event such as a landslide. The long life-span and slow and regular growth rate of some species can be used to date events (lichenometry). Lichens are a keystone species in many ecosystems and benefit trees and birds.
Etymology and pronunciation
The English word lichen derives from the Greek ("tree moss, lichen, lichen-like eruption on skin") via Latin . The Greek noun, which literally means "licker", derives from the verb , "to lick". In American English, "lichen" is pronounced the same as the verb "liken" (). In British English, both this pronunciation and one rhyming with "kitchen" () are used.
Anatomy and morphology
Growth forms
Lichens grow in a wide range of shapes and forms; this external appearance is known as their morphology. The shape of a lichen is usually determined by the organization of the fungal filaments. The nonreproductive tissues, or vegetative body parts, are called the thallus. Lichens are grouped by thallus type, since the thallus is usually the most visually prominent part of the lichen. Thallus growth forms typically correspond to a few basic internal structure types. Common names for lichens often come from a growth form or color that is typical of a lichen genus.
Common groupings of lichen thallus growth forms are:
fruticose – growing like a tuft or multiple-branched leafless mini-shrub, upright or hanging down, 3-dimensional branches with nearly round cross section (terete) or flattened
foliose – growing in 2-dimensional, flat, leaf-like lobes
crustose – crust-like, adhering tightly to a surface (substrate) like a thick coat of paint
squamulose – formed of small leaf-like scales crustose below but free at the tips
leprose – powdery
gelatinous – jelly-like
filamentous – stringy or like matted hair
byssoid – wispy, like teased wool
structureless
There are variations in growth types in a single lichen species, grey areas between the growth type descriptions, and overlapping between growth types, so some authors might describe lichens using different growth type descriptions.
When a crustose lichen gets old, the center may start to crack up like old-dried paint, old-broken asphalt paving, or like the polygonal "islands" of cracked-up mud in a dried lakebed. This is called being rimose or areolate, and the "island" pieces separated by the cracks are called areolas. The areolas appear separated, but are (or were) connected by an underlying prothallus or hypothallus. When a crustose lichen grows from a center and appears to radiate out, it is called crustose placodioid. When the edges of the areolas lift up from the substrate, it is called squamulose.
These growth form groups are not precisely defined. Foliose lichens may sometimes branch and appear to be fruticose. Fruticose lichens may have flattened branching parts and appear leafy. Squamulose lichens may appear where the edges lift up. Gelatinous lichens may appear leafy when dry.
The thallus is not always the part of the lichen that is most visually noticeable. Some lichens can grow inside solid rock between the grains (endolithic lichens), with only the sexual fruiting part visible growing outside the rock. These may be dramatic in color or appearance. Forms of these sexual parts are not in the above growth form categories. The most visually noticeable reproductive parts are often circular, raised, plate-like or disc-like outgrowths, with crinkly edges, and are described in sections below.
Color
Lichens come in many colors. Coloration is usually determined by the photosynthetic component. Special pigments, such as yellow usnic acid, give lichens a variety of colors, including reds, oranges, yellows, and browns, especially in exposed, dry habitats. In the absence of special pigments, lichens are usually bright green to olive gray when wet, gray or grayish-green to brown when dry. This is because moisture causes the surface skin (cortex) to become more transparent, exposing the green photobiont layer. Different colored lichens covering large areas of exposed rock surfaces, or lichens covering or hanging from bark can be a spectacular display when the patches of diverse colors "come to life" or "glow" in brilliant displays following rain.
Different colored lichens may inhabit different adjacent sections of a rock face, depending on the angle of exposure to light. Colonies of lichens may be spectacular in appearance, dominating much of the surface of the visual landscape in forests and natural places, such as the vertical "paint" covering the vast rock faces of Yosemite National Park.
Color is used in identification. The color of a lichen changes depending on whether the lichen is wet or dry. Color descriptions used for identification are based on the color that shows when the lichen is dry. Dry lichens with a cyanobacterium as the photosynthetic partner tend to be dark grey, brown, or black.
The underside of the leaf-like lobes of foliose lichens is a different color from the top side (dorsiventral), often brown or black, sometimes white. A fruticose lichen may have flattened "branches", appearing similar to a foliose lichen, but the underside of a leaf-like structure on a fruticose lichen is the same color as the top side. The leaf-like lobes of a foliose lichen may branch, giving the appearance of a fruticose lichen, but the underside will be a different color from the top side.
The sheen on some jelly-like gelatinous lichens is created by mucilaginous secretions.
Internal structure
A lichen consists of a simple photosynthesizing organism, usually a green alga or cyanobacterium, surrounded by filaments of a fungus. Generally, most of a lichen's bulk is made of interwoven fungal filaments, but this is reversed in filamentous and gelatinous lichens. The fungus is called a mycobiont. The photosynthesizing organism is called a photobiont. Algal photobionts are called phycobionts. Cyanobacteria photobionts are called cyanobionts.
The part of a lichen that is not involved in reproduction, the "body" or "vegetative tissue" of a lichen, is called the thallus. The thallus form is very different from any form where the fungus or alga are growing separately. The thallus is made up of filaments of the fungus called hyphae. The filaments grow by branching then rejoining to create a mesh, which is called being "anastomosed". The mesh of fungal filaments may be dense or loose.
Generally, the fungal mesh surrounds the algal or cyanobacterial cells, often enclosing them within complex fungal tissues that are unique to lichen associations. The thallus may or may not have a protective "skin" of densely packed fungal filaments, often containing a second fungal species, which is called a cortex. Fruticose lichens have one cortex layer wrapping around the "branches". Foliose lichens have an upper cortex on the top side of the "leaf", and a separate lower cortex on the bottom side. Crustose and squamulose lichens have only an upper cortex, with the "inside" of the lichen in direct contact with the surface they grow on (the substrate). Even if the edges peel up from the substrate and appear flat and leaf-like, they lack a lower cortex, unlike foliose lichens. Filamentous, byssoid, leprose, gelatinous, and other lichens do not have a cortex; in other words, they are ecorticate.
Fruticose, foliose, crustose, and squamulose lichens generally have up to three different types of tissue, differentiated by having different densities of fungal filaments. The top layer, where the lichen contacts the environment, is called a cortex. The cortex is made of densely tightly woven, packed, and glued together (agglutinated) fungal filaments. The dense packing makes the cortex act like a protective "skin", keeping other organisms out, and reducing the intensity of sunlight on the layers below. The cortex layer can be up to several hundred micrometers (μm) in thickness (less than a millimeter). The cortex may be further topped by an epicortex of secretions, not cells, 0.6–1 μm thick in some lichens. This secretion layer may or may not have pores.
Below the cortex layer is a layer called the photobiontic layer or symbiont layer. The symbiont layer has less densely packed fungal filaments, with the photosynthetic partner embedded in them. The less dense packing allows air circulation during photosynthesis, similar to the anatomy of a leaf. Each cell or group of cells of the photobiont is usually individually wrapped by hyphae, and in some cases penetrated by a haustorium. In crustose and foliose lichens, algae in the photobiontic layer are diffuse among the fungal filaments, decreasing in gradation into the layer below. In fruticose lichens, the photobiontic layer is sharply distinct from the layer below.
The layer beneath the symbiont layer is called the medulla. The medulla is less densely packed with fungal filaments than the layers above. In foliose lichens, as in Peltigera, there is usually another densely packed layer of fungal filaments called the lower cortex. Root-like fungal structures called rhizines (usually) grow from the lower cortex to attach or anchor the lichen to the substrate. Fruticose lichens have a single cortex wrapping all the way around the "stems" and "branches". The medulla is the lowest layer, and may form a cottony white inner core for the branchlike thallus, or it may be hollow. Crustose and squamulose lichens lack a lower cortex, and the medulla is in direct contact with the substrate that the lichen grows on.
In crustose areolate lichens, the edges of the areolas peel up from the substrate and appear leafy. In squamulose lichens the part of the lichen thallus that is not attached to the substrate may also appear leafy. But these leafy parts lack a lower cortex, which distinguishes crustose and squamulose lichens from foliose lichens. Conversely, foliose lichens may appear flattened against the substrate like a crustose lichen, but most of the leaf-like lobes can be lifted up from the substrate because it is separated from it by a tightly packed lower cortex.
Gelatinous, byssoid, and leprose lichens lack a cortex (are ecorticate), and generally have only undifferentiated tissue, similar to only having a symbiont layer.
In lichens that include both green algal and cyanobacterial symbionts, the cyanobacteria may be held on the upper or lower surface in small pustules called cephalodia.
Pruinia is a whitish coating on top of an upper surface. An epinecral layer is "a layer of horny dead fungal hyphae with indistinct lumina in or near the cortex above the algal layer".
In August 2016, it was reported that some macrolichens have more than one species of fungus in their tissues.
Physiology
Symbiotic relation
A lichen is a composite organism that emerges from algae or cyanobacteria living among the filaments (hyphae) of the fungi in a mutually beneficial symbiotic relationship. The fungi benefit from the carbohydrates produced by the algae or cyanobacteria via photosynthesis. The algae or cyanobacteria benefit by being protected from the environment by the filaments of the fungi, which also gather moisture and nutrients from the environment, and (usually) provide an anchor to it. Although some photosynthetic partners in a lichen can survive outside the lichen, the lichen symbiotic association extends the ecological range of both partners, whereby most descriptions of lichen associations describe them as symbiotic. Both partners gain water and mineral nutrients mainly from the atmosphere, through rain and dust. The fungal partner protects the alga by retaining water, serving as a larger capture area for mineral nutrients and, in some cases, provides minerals obtained from the substrate. If a cyanobacterium is present, as a primary partner or another symbiont in addition to a green alga as in certain tripartite lichens, they can fix atmospheric nitrogen, complementing the activities of the green alga.
In three different lineages the fungal partner has independently lost the mitochondrial gene atp9, which has key functions in mitochondrial energy production. The loss makes the fungi completely dependent on their symbionts.
The algal or cyanobacterial cells are photosynthetic and, as in plants, they reduce atmospheric carbon dioxide into organic carbon sugars to feed both symbionts. Phycobionts (algae) produce sugar alcohols (ribitol, sorbitol, and erythritol), which are absorbed by the mycobiont (fungus). Cyanobionts produce glucose. Lichenized fungal cells can make the photobiont "leak" out the products of photosynthesis, where they can then be absorbed by the fungus.
It appears many, probably the majority, of lichen also live in a symbiotic relationship with an order of basidiomycete yeasts called Cyphobasidiales. The absence of this third partner could explain why growing lichen in the laboratory is difficult. The yeast cells are responsible for the formation of the characteristic cortex of the lichen thallus, and could also be important for its shape.
The lichen combination of alga or cyanobacterium with a fungus has a very different form (morphology), physiology, and biochemistry than the component fungus, alga, or cyanobacterium growing by itself, naturally or in culture. The body (thallus) of most lichens is different from those of either the fungus or alga growing separately. When grown in the laboratory in the absence of its photobiont, a lichen fungus develops as a structureless, undifferentiated mass of fungal filaments (hyphae). If combined with its photobiont under appropriate conditions, its characteristic form associated with the photobiont emerges, in the process called morphogenesis. In a few remarkable cases, a single lichen fungus can develop into two very different lichen forms when associating with either a green algal or a cyanobacterial symbiont. Quite naturally, these alternative forms were at first considered to be different species, until they were found growing in a conjoined manner.
Evidence that lichens are examples of successful symbiosis is the fact that lichens can be found in almost every habitat and geographic area on the planet. Two species in two genera of green algae are found in over 35% of all lichens, but can only rarely be found living on their own outside of a lichen.
In a case where one fungal partner simultaneously had two green algae partners that outperform each other in different climates, this might indicate having more than one photosynthetic partner at the same time might enable the lichen to exist in a wider range of habitats and geographic locations.
At least one form of lichen, the North American beard-like lichens, are constituted of not two but three symbiotic partners: an ascomycetous fungus, a photosynthetic alga, and, unexpectedly, a basidiomycetous yeast.
Phycobionts can have a net output of sugars with only water vapor. The thallus must be saturated with liquid water for cyanobionts to photosynthesize.
Algae produce sugars that are absorbed by the fungus by diffusion into special fungal hyphae called appressoria or haustoria in contact with the wall of the algal cells. The appressoria or haustoria may produce a substance that increases permeability of the algal cell walls, and may penetrate the walls. The algae may contribute up to 80% of their sugar production to the fungus.
Ecology
Lichen associations may be examples of mutualism or commensalism, but the lichen relationship can be considered parasitic under circumstances where the photosynthetic partner can exist in nature independently of the fungal partner, but not vice versa. Photobiont cells are routinely destroyed in the course of nutrient exchange. The association continues because reproduction of the photobiont cells matches the rate at which they are destroyed. The fungus surrounds the algal cells, often enclosing them within complex fungal tissues unique to lichen associations. In many species the fungus penetrates the algal cell wall, forming penetration pegs (haustoria) similar to those produced by pathogenic fungi that feed on a host. Cyanobacteria in laboratory settings can grow faster when they are alone rather than when they are part of a lichen.
Miniature ecosystem and holobiont theory
Symbiosis in lichens is so well-balanced that lichens have been considered to be relatively self-contained miniature ecosystems in and of themselves. It is thought that lichens may be even more complex symbiotic systems that include non-photosynthetic bacterial communities performing other functions as partners in a holobiont.
Many lichens are very sensitive to environmental disturbances and can be used to cheaply assess air pollution, ozone depletion, and metal contamination. Lichens have been used in making dyes, perfumes (oakmoss), and in traditional medicines. A few lichen species are eaten by insects or larger animals, such as reindeer. Lichens are widely used as environmental indicators or bio-indicators. When air is very badly polluted with sulphur dioxide, there may be no lichens present; only some green algae can tolerate those conditions. If the air is clean, then shrubby, hairy and leafy lichens become abundant. A few lichen species can tolerate fairly high levels of pollution, and are commonly found in urban areas, on pavements, walls and tree bark. The most sensitive lichens are shrubby and leafy, while the most tolerant lichens are all crusty in appearance. Since industrialisation, many of the shrubby and leafy lichens such as Ramalina, Usnea and Lobaria species have very limited ranges, often being confined to the areas which have the cleanest air.
Lichenicolous fungi
Some fungi can only be found living on lichens as obligate parasites. These are referred to as lichenicolous fungi, and are a different species from the fungus living inside the lichen; thus they are not considered to be part of the lichen.
Reaction to water
Moisture makes the cortex become more transparent. This way, the algae can conduct photosynthesis when moisture is available, and is protected at other times. When the cortex is more transparent, the algae show more clearly and the lichen looks greener.
Metabolites, metabolite structures and bioactivity
Lichens can show intense antioxidant activity. Secondary metabolites are often deposited as crystals in the apoplast. Secondary metabolites are thought to play a role in preference for some substrates over others.
Growth rate
Lichens often have a regular but very slow growth rate of less than a millimeter per year.
In crustose lichens, the area along the margin is where the most active growth is taking place. Most crustose lichens grow only 1–2 mm in diameter per year.
Life span
Lichens may be long-lived, with some considered to be among the oldest living organisms. Lifespan is difficult to measure because what defines the "same" individual lichen is not precise. Lichens grow by vegetatively breaking off a piece, which may or may not be defined as the "same" lichen, and two lichens can merge, then becoming the "same" lichen. One specimen of Rhizocarpon geographicum on East Baffin Island has an estimated age of 9500 years. Thalli of Rhizocarpon geographicum and Rhizocarpon eupetraeoides/inarense in the central Brooks Range of northern Alaska have been given a maximum possible age of 10,000–11,500 years.
Response to environmental stress
Unlike simple dehydration in plants and animals, lichens may experience a complete loss of body water in dry periods. Lichens are capable of surviving extremely low levels of water content (poikilohydric). They quickly absorb water when it becomes available again, becoming soft and fleshy.
In tests, lichen survived and showed remarkable results on the adaptation capacity of photosynthetic activity within the simulation time of 34 days under Martian conditions in the Mars Simulation Laboratory (MSL) maintained by the German Aerospace Center (DLR).
The European Space Agency has discovered that lichens can survive unprotected in space. In an experiment led by Leopoldo Sancho from the Complutense University of Madrid, two species of lichen—Rhizocarpon geographicum and Rusavskia elegans—were sealed in a capsule and launched on a Russian Soyuz rocket 31 May 2005. Once in orbit, the capsules were opened and the lichens were directly exposed to the vacuum of space with its widely fluctuating temperatures and cosmic radiation. After 15 days, the lichens were brought back to earth and were found to be unchanged in their ability to photosynthesize.
Reproduction and dispersal
Vegetative reproduction
Many lichens reproduce asexually, either by a piece breaking off and growing on its own (vegetative reproduction) or through the dispersal of diaspores containing a few algal cells surrounded by fungal cells. Because of the relative lack of differentiation in the thallus, the line between diaspore formation and vegetative reproduction is often blurred. Fruticose lichens can fragment, and new lichens can grow from the fragment (vegetative reproduction). Many lichens break up into fragments when they dry, dispersing themselves by wind action, to resume growth when moisture returns. Soredia (singular: "soredium") are small groups of algal cells surrounded by fungal filaments that form in structures called soralia, from which the soredia can be dispersed by wind. Isidia (singular: "isidium") are branched, spiny, elongated, outgrowths from the thallus that break off for mechanical dispersal. Lichen propagules (diaspores) typically contain cells from both partners, although the fungal components of so-called "fringe species" rely instead on algal cells dispersed by the "core species".
Sexual reproduction
Structures involved in reproduction often appear as discs, bumps, or squiggly lines on the surface of the thallus. Though it has been argued that sexual reproduction in photobionts is selected against, there is strong evidence that suggests meiotic activities (sexual reproduction) in Trebouxia. Many lichen fungi reproduce sexually like other fungi, producing spores formed by meiosis and fusion of gametes. Following dispersal, such fungal spores must meet with a compatible algal partner before a functional lichen can form.
Some lichen fungi belong to the phylum Basidiomycota (basidiolichens) and produce mushroom-like reproductive structures resembling those of their nonlichenized relatives.
Most lichen fungi belong to Ascomycetes (ascolichens). Among the ascolichens, spores are produced in spore-producing structures called ascomata. The most common types of ascomata are the apothecium (plural: apothecia) and perithecium (plural: perithecia). Apothecia are usually cups or plate-like discs located on the top surface of the lichen thallus. When apothecia are shaped like squiggly line segments instead of like discs, they are called lirellae. Perithecia are shaped like flasks that are immersed in the lichen thallus tissue, which has a small hole for the spores to escape the flask, and appear like black dots on the lichen surface.
The three most common spore body types are raised discs called apothecia (singular: apothecium), bottle-like cups with a small hole at the top called perithecia (singular: perithecium), and pycnidia (singular: pycnidium), shaped like perithecia but without asci (an ascus is the structure that contains and releases the sexual spores in fungi of the Ascomycota).
The apothecium has a layer of exposed spore-producing cells called asci (singular: ascus), and is usually a different color from the thallus tissue. When the apothecium has an outer margin, the margin is called the exciple. When the exciple has a color similar to colored thallus tissue the apothecium or lichen is called lecanorine, meaning similar to members of the genus Lecanora. When the exciple is blackened like carbon it is called lecideine meaning similar to members of the genus Lecidea. When the margin is pale or colorless it is called biatorine.
A "podetium" (plural: podetia) is a lichenized stalk-like structure of the fruiting body rising from the thallus, associated with some fungi that produce a fungal apothecium. Since it is part of the reproductive tissue, podetia are not considered part of the main body (thallus), but may be visually prominent. The podetium may be branched, and sometimes cup-like. They usually bear the fungal pycnidia or apothecia or both. Many lichens have apothecia that are visible to the naked eye.
Most lichens produce abundant sexual structures. Many species appear to disperse only by sexual spores. For example, the crustose lichens Graphis scripta and Ochrolechia parella produce no symbiotic vegetative propagules. Instead, the lichen-forming fungi of these species reproduce sexually by self-fertilization (i.e. they are homothallic). This breeding system may enable successful reproduction in harsh environments.
Mazaedia (singular: mazaedium) are apothecia shaped like a dressmaker's pin in pin lichens, where the fruiting body is a brown or black mass of loose ascospores enclosed by a cup-shaped exciple, which sits on top of a tiny stalk.
Taxonomy and classification
Lichens are classified by the fungal component. Lichen species are given the same scientific name (binomial name) as the fungus species in the lichen. Lichens are being integrated into the classification schemes for fungi. The alga bears its own scientific name, which bears no relationship to that of the lichen or fungus. There are about 20,000 identified lichen species, and taxonomists have estimated that the total number of lichen species (including those yet undiscovered) might be as high as 28,000. Nearly 20% of known fungal species are associated with lichens.
"Lichenized fungus" may refer to the entire lichen, or to just the fungus. This may cause confusion without context. A particular fungus species may form lichens with different algae species, giving rise to what appear to be different lichen species, but which are still classified (as of 2014) as the same lichen species.
Formerly, some lichen taxonomists placed lichens in their own division, the Mycophycophyta, but this practice is no longer accepted because the components belong to separate lineages. Neither the ascolichens nor the basidiolichens form monophyletic lineages in their respective fungal phyla, but they do form several major solely or primarily lichen-forming groups within each phylum. Even more unusual than basidiolichens is the fungus Geosiphon pyriforme, a member of the Glomeromycota that is unique in that it encloses a cyanobacterial symbiont inside its cells. Geosiphon is not usually considered to be a lichen, and its peculiar symbiosis was not recognized for many years. The genus is more closely allied to endomycorrhizal genera. Fungi from Verrucariales also form marine lichens with the brown algae Petroderma maculiforme, and have a symbiotic relationship with seaweed (such as rockweed) and Blidingia minima, where the algae are the dominant components. The fungi is thought to help the rockweeds to resist desiccation when exposed to air. In addition, lichens can also use yellow-green algae (Heterococcus) as their symbiotic partner.
Lichens independently emerged from fungi associating with algae and cyanobacteria multiple times throughout history.
Fungi
The fungal component of a lichen is called the mycobiont. The mycobiont may be an Ascomycete or Basidiomycete. The associated lichens are called either ascolichens or basidiolichens, respectively. Living as a symbiont in a lichen appears to be a successful way for a fungus to derive essential nutrients, since about 20% of all fungal species have acquired this mode of life.
Thalli produced by a given fungal symbiont with its differing partners may be similar, and the secondary metabolites identical, indicating that the fungus has the dominant role in determining the morphology of the lichen. But the same mycobiont with different photobionts may also produce very different growth forms. Lichens are known in which there is one fungus associated with two or even three algal species.
Although each lichen thallus generally appears homogeneous, some evidence seems to suggest that the fungal component may consist of more than one genetic individual of that species.
Two or more fungal species can interact to form the same lichen.
The following table lists the orders and families of fungi that include lichen-forming species.
Photobionts
The photosynthetic partner in a lichen is called a photobiont. The photobionts in lichens come from a variety of simple prokaryotic and eukaryotic organisms. In the majority of lichens the photobiont is a green alga (Chlorophyta) or a cyanobacterium. In some lichens both types are present; in such cases, the alga is typically the primary partner, with the cyanobacteria being located in cryptic pockets. Algal photobionts are called phycobionts, while cyanobacterial photobionts are called cyanobionts. About 90% of all known lichens have phycobionts, and about 10% have cyanobionts. Approximately 100 species of photosynthetic partners from 40 genera and five distinct classes (prokaryotic: Cyanophyceae; eukaryotic: Trebouxiophyceae, Phaeophyceae, Chlorophyceae) have been found to associate with the lichen-forming fungi.
Common algal photobionts are from the genera Trebouxia, Trentepohlia, Pseudotrebouxia, or Myrmecia. Trebouxia is the most common genus of green algae in lichens, occurring in about 40% of all lichens. "Trebouxioid" means either a photobiont that is in the genus Trebouxia, or resembles a member of that genus, and is therefore presumably a member of the class Trebouxiophyceae. The second most commonly represented green alga genus is Trentepohlia. Overall, about 100 species of eukaryotes are known to occur as photobionts in lichens. All the algae are probably able to exist independently in nature as well as in the lichen.
A "cyanolichen" is a lichen with a cyanobacterium as its main photosynthetic component (photobiont). Most cyanolichen are also ascolichens, but a few basidiolichen like Dictyonema and Acantholichen have cyanobacteria as their partner.
The most commonly occurring cyanobacterium genus is Nostoc. Other common cyanobacterium photobionts are from Scytonema. Many cyanolichens are small and black, and have limestone as the substrate. Another cyanolichen group, the jelly lichens of the genera Collema or Leptogium are gelatinous and live on moist soils. Another group of large and foliose species including Peltigera, Lobaria, and Degelia are grey-blue, especially when dampened or wet. Many of these characterize the Lobarion communities of higher rainfall areas in western Britain, e.g., in the Celtic rain forest. Strains of cyanobacteria found in various cyanolichens are often closely related to one another. They differ from the most closely related free-living strains.
The lichen association is a close symbiosis. It extends the ecological range of both partners but is not always obligatory for their growth and reproduction in natural environments, since many of the algal symbionts can live independently. A prominent example is the alga Trentepohlia, which forms orange-coloured populations on tree trunks and suitable rock faces. Lichen propagules (diaspores) typically contain cells from both partners, although the fungal components of so-called "fringe species" rely instead on algal cells dispersed by the "core species".
The same cyanobiont species can occur in association with different fungal species as lichen partners. The same phycobiont species can occur in association with different fungal species as lichen partners. More than one phycobiont may be present in a single thallus.
A single lichen may contain several algal genotypes. These multiple genotypes may better enable response to adaptation to environmental changes, and enable the lichen to inhabit a wider range of environments.
Controversy over classification method and species names
There are about 20,000 known lichen species. But what is meant by "species" is different from what is meant by biological species in plants, animals, or fungi, where being the same species implies that there is a common ancestral lineage. Because lichens are combinations of members of two or even three different biological kingdoms, these components must have a different ancestral lineage from each other. By convention, lichens are still called "species" anyway, and are classified according to the species of their fungus, not the species of the algae or cyanobacteria. Lichens are given the same scientific name (binomial name) as the fungus in them, which may cause some confusion. The alga bears its own scientific name, which has no relationship to the name of the lichen or fungus.
Depending on context, "lichenized fungus" may refer to the entire lichen, or to the fungus when it is in the lichen, which can be grown in culture in isolation from the algae or cyanobacteria. Some algae and cyanobacteria are found naturally living outside of the lichen. The fungal, algal, or cyanobacterial component of a lichen can be grown by itself in culture. When growing by themselves, the fungus, algae, or cyanobacteria have very different properties than those of the lichen. Lichen properties such as growth form, physiology, and biochemistry, are very different from the combination of the properties of the fungus and the algae or cyanobacteria.
The same fungus growing in combination with different algae or cyanobacteria, can produce lichens that are very different in most properties, meeting non-DNA criteria for being different "species". Historically, these different combinations were classified as different species. When the fungus is identified as being the same using modern DNA methods, these apparently different species get reclassified as the same species under the current (2014) convention for classification by fungal component. This has led to debate about this classification convention. These apparently different "species" have their own independent evolutionary history.
There is also debate as to the appropriateness of giving the same binomial name to the fungus, and to the lichen that combines that fungus with an alga or cyanobacterium (synecdoche). This is especially the case when combining the same fungus with different algae or cyanobacteria produces dramatically different lichen organisms, which would be considered different species by any measure other than the DNA of the fungal component. If the whole lichen produced by the same fungus growing in association with different algae or cyanobacteria, were to be classified as different "species", the number of "lichen species" would be greater.
Diversity
The largest number of lichenized fungi occur in the Ascomycota, with about 40% of species forming such an association. Some of these lichenized fungi occur in orders with nonlichenized fungi that live as saprotrophs or plant parasites (for example, the Leotiales, Dothideales, and Pezizales). Other lichen fungi occur in only five orders in which all members are engaged in this habit (Orders Graphidales, Gyalectales, Peltigerales, Pertusariales, and Teloschistales). Overall, about 98% of lichens have an ascomycetous mycobiont. Next to the Ascomycota, the largest number of lichenized fungi occur in the unassigned fungi imperfecti, a catch-all category for fungi whose sexual form of reproduction has never been observed. Comparatively few basidiomycetes are lichenized, but these include agarics, such as species of Lichenomphalia, clavarioid fungi, such as species of Multiclavula, and corticioid fungi, such as species of Dictyonema.
Identification methods
Lichen identification uses growth form, microscopy and reactions to chemical tests.
The outcome of the "Pd test" is called "Pd", which is also used as an abbreviation for the chemical used in the test, para-phenylenediamine. If putting a drop on a lichen turns an area bright yellow to orange, this helps identify it as belonging to either the genus Cladonia or Lecanora.
Evolution and paleontology
The fossil record for lichens is poor. The extreme habitats that lichens dominate, such as tundra, mountains, and deserts, are not ordinarily conducive to producing fossils. There are fossilized lichens embedded in amber. The fossilized Anzia is found in pieces of amber in northern Europe and dates back approximately 40 million years. Lichen fragments are also found in fossil leaf beds, such as Lobaria from Trinity County in northern California, US, dating back to the early to middle Miocene.
The oldest fossil lichen in which both symbiotic partners have been recovered is Winfrenatia, an early zygomycetous (Glomeromycotan) lichen symbiosis that may have involved controlled parasitism, is permineralized in the Rhynie Chert of Scotland, dating from early Early Devonian, about 400 million years ago. The slightly older fossil Spongiophyton has also been interpreted as a lichen on morphological and isotopic grounds, although the isotopic basis is decidedly shaky. It has been demonstrated that Silurian-Devonian fossils Nematothallus and Prototaxites were lichenized. Thus lichenized Ascomycota and Basidiomycota were a component of Early Silurian-Devonian terrestrial ecosystems. Newer research suggests that lichen evolved after the evolution of land plants.
The ancestral ecological state of both Ascomycota and Basidiomycota was probably saprobism, and independent lichenization events may have occurred multiple times. In 1995, Gargas and colleagues proposed that there were at least five independent origins of lichenization; three in the basidiomycetes and at least two in the Ascomycetes. Lutzoni et al. (2001) suggest lichenization probably evolved earlier and was followed by multiple independent losses. Some non-lichen-forming fungi may have secondarily lost the ability to form a lichen association. As a result, lichenization has been viewed as a highly successful nutritional strategy.
Lichenized Glomeromycota may extend well back into the Precambrian. Lichen-like fossils consisting of coccoid cells (cyanobacteria?) and thin filaments (mucoromycotinan Glomeromycota?) are permineralized in marine phosphorite of the Doushantuo Formation in southern China. These fossils are thought to be 551 to 635 million years old or Ediacaran. Ediacaran acritarchs also have many similarities with Glomeromycotan vesicles and spores. It has also been claimed that Ediacaran fossils including Dickinsonia, were lichens, although this claim is controversial. Endosymbiotic Glomeromycota comparable with living Geosiphon may extend back into the Proterozoic in the form of 1500 million year old Horodyskia and 2200 million year old Diskagma. Discovery of these fossils suggest that fungi developed symbiotic partnerships with photoautotrophs long before the evolution of vascular plants, though the Ediacaran lichen hypothesis is largely rejected due to an inappropriate definition of lichens based on taphonomy and substrate ecology. However, a 2019 study by the same scientist who rejected the Ediacaran lichen hypothesis, Nelsen, used new time-calibrated phylogenies to conclude that there is no evidence of lichen before the existence of vascular plants.
Lecanoromycetes, one of the most common classes of lichen-forming fungi, diverged from its ancestor, which may have also been lichen forming, around 258 million years ago, during the late Paleozoic period. However, the closely related clade Euritiomycetes appears to have become lichen-forming only 52 million years ago, during the early Cenozoic period.
Ecology and interactions with environment
Substrates and habitats
Lichens grow on and in a wide range of substrates and habitats, including some of the most extreme conditions on earth. They are abundant growing on bark, leaves, and hanging from epiphyte branches in rain forests and in temperate woodland. They grow on bare rock, walls, gravestones, roofs, and exposed soil surfaces. They can survive in some of the most extreme environments on Earth: arctic tundra, hot dry deserts, rocky coasts, and toxic slag heaps. They can live inside solid rock, growing between the grains, and in the soil as part of a biological soil crust in arid habitats such as deserts. Some lichens do not grow on anything, living out their lives blowing about the environment.
When growing on mineral surfaces, some lichens slowly decompose their substrate by chemically degrading and physically disrupting the minerals, contributing to the process of weathering by which rocks are gradually turned into soil. While this contribution to weathering is usually benign, it can cause problems for artificial stone structures. For example, there is an ongoing lichen growth problem on Mount Rushmore National Memorial that requires the employment of mountain-climbing conservators to clean the monument.
Lichens are not parasites on the plants they grow on, but only use them as a substrate. The fungi of some lichen species may "take over" the algae of other lichen species. Lichens make their own food from their photosynthetic parts and by absorbing minerals from the environment. Lichens growing on leaves may have the appearance of being parasites on the leaves, but they are not. Some lichens in Diploschistes parasitise other lichens. Diploschistes muscorum starts its development in the tissue of a host Cladonia species.
In the arctic tundra, lichens, together with mosses and liverworts, make up the majority of the ground cover, which helps insulate the ground and may provide forage for grazing animals. An example is "reindeer moss", which is a lichen, not a moss.
There are only two species of known permanently submerged lichens; Hydrothyria venosa is found in fresh water environments, and Verrucaria serpuloides is found in marine environments.
A crustose lichen that grows on rock is called a saxicolous lichen. Crustose lichens that grow on the rock are epilithic, and those that grow immersed inside rock, growing between the crystals with only their fruiting bodies exposed to the air, are called endolithic lichens. A crustose lichen that grows on bark is called a corticolous lichen. A lichen that grows on wood from which the bark has been stripped is called a lignicolous lichen. Lichens that grow immersed inside plant tissues are called endophloidic lichens or endophloidal lichens. Lichens that use leaves as substrates, whether the leaf is still on the tree or on the ground, are called epiphyllous or foliicolous. A terricolous lichen grows on the soil as a substrate. Many squamulose lichens are terricolous. Umbilicate lichens are foliose lichens that are attached to the substrate at only one point. A vagrant lichen is not attached to a substrate at all, and lives its life being blown around by the wind.
Lichens and soils
In addition to distinct physical mechanisms by which lichens break down raw stone, studies indicate lichens attack stone chemically, entering newly chelated minerals into the ecology. The substances exuded by lichens, known for their strong ability to bind and sequester metals, along with the common formation of new minerals, especially metal oxalates, and the traits of the substrates they alter, all highlight the important role lichens play in the process of chemical weathering. Over time, this activity creates new fertile soil from stone.
Lichens may be important in contributing nitrogen to soils in some deserts through being eaten, along with their rock substrate, by snails, which then defecate, putting the nitrogen into the soils. Lichens help bind and stabilize soil sand in dunes. In deserts and semi-arid areas, lichens are part of extensive, living biological soil crusts, essential for maintaining the soil structure.
Ecological interactions
Lichens are pioneer species, among the first living things to grow on bare rock or areas denuded of life by a disaster. Lichens may have to compete with plants for access to sunlight, but because of their small size and slow growth, they thrive in places where higher plants have difficulty growing. Lichens are often the first to settle in places lacking soil, constituting the sole vegetation in some extreme environments such as those found at high mountain elevations and at high latitudes. Some survive in the tough conditions of deserts, and others on frozen soil of the Arctic regions.
A major ecophysiological advantage of lichens is that they are poikilohydric (poikilo- variable, hydric- relating to water), meaning that though they have little control over the status of their hydration, they can tolerate irregular and extended periods of severe desiccation. Like some mosses, liverworts, ferns and a few resurrection plants, upon desiccation, lichens enter a metabolic suspension or stasis (known as cryptobiosis) in which the cells of the lichen symbionts are dehydrated to a degree that halts most biochemical activity. In this cryptobiotic state, lichens can survive wider extremes of temperature, radiation and drought in the harsh environments they often inhabit.
Lichens do not have roots and do not need to tap continuous reservoirs of water like most higher plants, thus they can grow in locations impossible for most plants, such as bare rock, sterile soil or sand, and various artificial structures such as walls, roofs, and monuments. Many lichens also grow as epiphytes (epi- on the surface, phyte- plant) on plants, particularly on the trunks and branches of trees. When growing on plants, lichens are not parasites; they do not consume any part of the plant nor poison it. Lichens produce allelopathic chemicals that inhibit the growth of mosses. Some ground-dwelling lichens, such as members of the subgenus Cladina (reindeer lichens), produce allelopathic chemicals that leach into the soil and inhibit the germination of seeds, spruce and other plants. Stability (that is, longevity) of their substrate is a major factor of lichen habitats. Most lichens grow on stable rock surfaces or the bark of old trees, but many others grow on soil and sand. In these latter cases, lichens are often an important part of soil stabilization; indeed, in some desert ecosystems, vascular (higher) plant seeds cannot become established except in places where lichen crusts stabilize the sand and help retain water.
Lichens may be eaten by some animals, such as reindeer, living in arctic regions. The larvae of a number of Lepidoptera species feed exclusively on lichens. These include common footman and marbled beauty. They are very low in protein and high in carbohydrates, making them unsuitable for some animals. The Northern flying squirrel uses it for nesting, food and winter water.
Effects of air pollution
If lichens are exposed to air pollutants at all times, without any deciduous parts, they are unable to avoid the accumulation of pollutants. Also lacking stomata and a cuticle, lichens may absorb aerosols and gases over the entire thallus surface from which they may readily diffuse to the photobiont layer. Because lichens do not possess roots, their primary source of most elements is the air, and therefore elemental levels in lichens often reflect the accumulated composition of ambient air. The processes by which atmospheric deposition occurs include fog and dew, gaseous absorption, and dry deposition. Consequently, environmental studies with lichens emphasize their feasibility as effective biomonitors of atmospheric quality.
Not all lichens are equally sensitive to air pollutants, so different lichen species show different levels of sensitivity to specific atmospheric pollutants. The sensitivity of a lichen to air pollution is directly related to the energy needs of the mycobiont, so that the stronger the dependency of the mycobiont on the photobiont, the more sensitive the lichen is to air pollution. Upon exposure to air pollution, the photobiont may use metabolic energy for repair of its cellular structures that would otherwise be used for maintenance of its photosynthetic activity, therefore leaving less metabolic energy available for the mycobiont. The alteration of the balance between the photobiont and mycobiont can lead to the breakdown of the symbiotic association. Therefore, lichen decline may result not only from the accumulation of toxic substances, but also from altered nutrient supplies that favor one symbiont over the other.
This interaction between lichens and air pollution has been used as a means of monitoring air quality since 1859, with more systematic methods developed by William Nylander in 1866.
Human use
Food
Lichens are eaten by many different cultures across the world. Although some lichens are only eaten in times of famine, others are a staple food or even a delicacy. Two obstacles are often encountered when eating lichens: lichen polysaccharides are generally indigestible to humans, and lichens usually contain mildly toxic secondary compounds that should be removed before eating. Very few lichens are poisonous, but those high in vulpinic acid or usnic acid are toxic. Most poisonous lichens are yellow.
In the past, Iceland moss (Cetraria islandica) was an important source of food for humans in northern Europe, and was cooked as a bread, porridge, pudding, soup, or salad. Bryoria fremontii (edible horsehair lichen) was an important food in parts of North America, where it was usually pitcooked. Northern peoples in North America and Siberia traditionally eat the partially digested reindeer lichen (Cladina spp.) after they remove it from the rumen of caribou or reindeer that have been killed. Rock tripe (Umbilicaria spp. and Lasalia spp.) is a lichen that has frequently been used as an emergency food in North America, and one species, Umbilicaria esculenta, (iwatake in Japanese) is used in a variety of traditional Korean and Japanese foods.
Lichenometry
Lichenometry is a technique used to determine the age of exposed rock surfaces based on the size of lichen thalli. Introduced by Beschel in the 1950s, the technique has found many applications. it is used in archaeology, palaeontology, and geomorphology. It uses the presumed regular but slow rate of lichen growth to determine the age of exposed rock. Measuring the diameter (or other size measurement) of the largest lichen of a species on a rock surface indicates the length of time since the rock surface was first exposed. Lichen can be preserved on old rock faces for up to 10,000 years, providing the maximum age limit of the technique, though it is most accurate (within 10% error) when applied to surfaces that have been exposed for less than 1,000 years. Lichenometry is especially useful for dating surfaces less than 500 years old, as radiocarbon dating techniques are less accurate over this period. The lichens most commonly used for lichenometry are those of the genera Rhizocarpon (e.g. the species Rhizocarpon geographicum, map lichen) and Xanthoria.
Biodegradation
Lichens have been shown to degrade polyester resins, as can be seen in archaeological sites in the Roman city of Baelo Claudia in Spain. Lichens can accumulate several environmental pollutants such as lead, copper, and radionuclides. Some species of lichen, such as Parmelia sulcata (called a hammered shield lichen, among other names) and Lobaria pulmonaria (lung lichen), and many in the Cladonia genus, have been shown to produce serine proteases capable of the degradation of pathogenic forms of prion protein (PrP), which may be useful in treating contaminated environmental reservoirs.
Dyes
Many lichens produce secondary compounds, including pigments that reduce harmful amounts of sunlight and powerful toxins that deter herbivores or kill bacteria. These compounds are very useful for lichen identification, and have had economic importance as dyes such as cudbear or primitive antibiotics.
A pH indicator (which can indicate acidic or basic substances) called litmus is a dye extracted from the lichen Roccella tinctoria ("dyer's weed") by boiling. It gives its name to the well-known litmus test.
Traditional dyes of the Scottish Highlands for Harris tweed and other traditional cloths were made from lichens, including the orange Xanthoria parietina ("common orange lichen") and the grey foliaceous Parmelia saxatilis common on rocks and known colloquially as "crottle".
There are reports dating almost 2,000 years old of lichens being used to make purple and red dyes. Of great historical and commercial significance are lichens belonging to the family Roccellaceae, commonly called orchella weed or orchil. Orcein and other lichen dyes have largely been replaced by synthetic versions.
Traditional medicine and research
Historically, in traditional medicine of Europe, Lobaria pulmonaria was collected in large quantities as "lungwort", due to its lung-like appearance (the "doctrine of signatures" suggesting that herbs can treat body parts that they physically resemble).Similarly, Peltigera leucophlebia ("ruffled freckled pelt") was used as a supposed cure for thrush, due to the resemblance of its cephalodia to the appearance of the disease.
Lichens produce metabolites being researched for their potential therapeutic or diagnostic value. Some metabolites produced by lichens are structurally and functionally similar to broad-spectrum antibiotics while few are associated respectively to antiseptic similarities. Usnic acid is the most commonly studied metabolite produced by lichens. It is also under research as a bactericidal agent against Escherichia coli and Staphylococcus aureus.
Aesthetic appeal
Colonies of lichens may be spectacular in appearance, dominating the surface of the visual landscape as part of the aesthetic appeal to visitors of Yosemite National Park, Sequoia National Park, and the Bay of Fires. Orange and yellow lichens add to the ambience of desert trees, tundras, and rocky seashores. Intricate webs of lichens hanging from tree branches add a mysterious aspect to forests. Fruticose lichens are used in model railroading and other modeling hobbies as a material for making miniature trees and shrubs.
In literature
In early Midrashic literature, the Hebrew word "vayilafeth" in Ruth 3:8 is explained as referring to Ruth entwining herself around Boaz like lichen. The 10th century Arab physician Al-Tamimi mentions lichens dissolved in vinegar and rose water being used in his day for the treatment of skin diseases and rashes.
The plot of John Wyndham's science fiction novel Trouble with Lichen revolves around an anti-aging chemical extracted from a lichen.
History
Although lichens had been recognized as organisms for quite some time, it was not until 1867, when Swiss botanist Simon Schwendener proposed his dual theory of lichens, that lichens are a combination of fungi with algae or cyanobacteria, whereby the true nature of the lichen association began to emerge. Schwendener's hypothesis, which at the time lacked experimental evidence, arose from his extensive analysis of the anatomy and development in lichens, algae, and fungi using a light microscope. Many of the leading lichenologists at the time, such as James Crombie and Nylander, rejected Schwendener's hypothesis because the consensus was that all living organisms were autonomous.
Other prominent biologists, such as Heinrich Anton de Bary, Albert Bernhard Frank, Beatrix Potter, Melchior Treub and Hermann Hellriegel, were not so quick to reject Schwendener's ideas and the concept soon spread into other areas of study, such as microbial, plant, animal and human pathogens. When the complex relationships between pathogenic microorganisms and their hosts were finally identified, Schwendener's hypothesis began to gain popularity. Further experimental proof of the dual nature of lichens was obtained when Eugen Thomas published his results in 1939 on the first successful re-synthesis experiment.
In the 2010s, a new facet of the fungi–algae partnership was discovered. Toby Spribille and colleagues found that many types of lichen that were long thought to be ascomycete–algae pairs were actually ascomycete–basidiomycete–algae trios. The third symbiotic partner in many lichens is a basidiomycete yeast.
| Biology and health sciences | Other organisms | null |
172474 | https://en.wikipedia.org/wiki/Lactic%20acid | Lactic acid | Lactic acid is an organic acid. It has the molecular formula C3H6O3. It is white in the solid state and it is miscible with water. When in the dissolved state, it forms a colorless solution. Production includes both artificial synthesis as well as natural sources. Lactic acid is an alpha-hydroxy acid (AHA) due to the presence of a hydroxyl group adjacent to the carboxyl group. It is used as a synthetic intermediate in many organic synthesis industries and in various biochemical industries. The conjugate base of lactic acid is called lactate (or the lactate anion). The name of the derived acyl group is lactoyl.
In solution, it can ionize by a loss of a proton to produce the lactate ion . Compared to acetic acid, its pK is 1 unit less, meaning lactic acid is ten times more acidic than acetic acid. This higher acidity is the consequence of the intramolecular hydrogen bonding between the α-hydroxyl and the carboxylate group.
Lactic acid is chiral, consisting of two enantiomers. One is known as -lactic acid, (S)-lactic acid, or (+)-lactic acid, and the other, its mirror image, is -lactic acid, (R)-lactic acid, or (−)-lactic acid. A mixture of the two in equal amounts is called -lactic acid, or racemic lactic acid. Lactic acid is hygroscopic. -Lactic acid is miscible with water and with ethanol above its melting point, which is about . -Lactic acid and -lactic acid have a higher melting point. Lactic acid produced by fermentation of milk is often racemic, although certain species of bacteria produce solely -lactic acid. On the other hand, lactic acid produced by fermentation in animal muscles has the () enantiomer and is sometimes called "sarcolactic" acid, from the Greek , meaning "flesh".
In animals, -lactate is constantly produced from pyruvate via the enzyme lactate dehydrogenase (LDH) in a process of fermentation during normal metabolism and exercise. It does not increase in concentration until the rate of lactate production exceeds the rate of lactate removal, which is governed by a number of factors, including monocarboxylate transporters, concentration and isoform of LDH, and oxidative capacity of tissues. The concentration of blood lactate is usually at rest, but can rise to over 20mM during intense exertion and as high as 25mM afterward. In addition to other biological roles, -lactic acid is the primary endogenous agonist of hydroxycarboxylic acid receptor 1 (HCA), which is a G protein-coupled receptor (GPCR).
In industry, lactic acid fermentation is performed by lactic acid bacteria, which convert simple carbohydrates such as glucose, sucrose, or galactose to lactic acid. These bacteria can also grow in the mouth; the acid they produce is responsible for the tooth decay known as cavities. In medicine, lactate is one of the main components of lactated Ringer's solution and Hartmann's solution. These intravenous fluids consist of sodium and potassium cations along with lactate and chloride anions in solution with distilled water, generally in concentrations isotonic with human blood. It is most commonly used for fluid resuscitation after blood loss due to trauma, surgery, or burns.
Lactic acid is produced in human tissues when the demand for oxygen is limited by the supply. This occurs during tissue ischemia when the flow of blood is limited as in sepsis or hemorrhagic shock. It may also occur when demand for oxygen is high such as with intense exercise. The process of lactic acidosis produces lactic acid which results in an oxygen debt which can be resolved or repaid when tissue oxygenation improves.
History
Swedish chemist Carl Wilhelm Scheele was the first person to isolate lactic acid in 1780 from sour milk. The name reflects the lact- combining form derived from the Latin word , meaning "milk". In 1808, Jöns Jacob Berzelius discovered that lactic acid (actually -lactate) is also produced in muscles during exertion. Its structure was established by Johannes Wislicenus in 1873.
In 1856, the role of Lactobacillus in the synthesis of lactic acid was discovered by Louis Pasteur. This pathway was used commercially by the German pharmacy Boehringer Ingelheim in 1895.
In 2006, global production of lactic acid reached 275,000 tonnes with an average annual growth of 10%.
Production
Lactic acid is produced industrially by bacterial fermentation of carbohydrates, or by chemical synthesis from acetaldehyde. , lactic acid was produced predominantly (70–90%) by fermentation. Production of racemic lactic acid consisting of a 1:1 mixture of and stereoisomers, or of mixtures with up to 99.9% -lactic acid, is possible by microbial fermentation. Industrial scale production of -lactic acid by fermentation is possible, but much more challenging.
Fermentative production
Fermented milk products are obtained industrially by fermentation of milk or whey by Lactobacillus bacteria: Lactobacillus acidophilus, Lacticaseibacillus casei (Lactobacillus casei), Lactobacillus delbrueckii subsp. bulgaricus (Lactobacillus bulgaricus), Lactobacillus helveticus, Lactococcus lactis , Bacillus amyloliquefaciens, and Streptococcus salivarius subsp. thermophilus (Streptococcus thermophilus).
As a starting material for industrial production of lactic acid, almost any carbohydrate source containing (Pentose sugar) and (Hexose sugar) can be used. Pure sucrose, glucose from starch, raw sugar, and beet juice are frequently used. Lactic acid producing bacteria can be divided in two classes: homofermentative bacteria like Lactobacillus casei and Lactococcus lactis, producing two moles of lactate from one mole of glucose, and heterofermentative species producing one mole of lactate from one mole of glucose as well as carbon dioxide and acetic acid/ethanol.
Chemical production
Racemic lactic acid is synthesized industrially by reacting acetaldehyde with hydrogen cyanide and hydrolysing the resultant lactonitrile. When hydrolysis is performed by hydrochloric acid, ammonium chloride forms as a by-product; the Japanese company Musashino is one of the last big manufacturers of lactic acid by this route. Synthesis of both racemic and enantiopure lactic acids is also possible from other starting materials (vinyl acetate, glycerol, etc.) by application of catalytic procedures.
Biology
Molecular biology
-Lactic acid is the primary endogenous agonist of hydroxycarboxylic acid receptor 1 (HCA1), a G protein-coupled receptor (GPCR).
Metabolism and exercise
During power exercises such as sprinting, when the rate of demand for energy is high, glucose is broken down and oxidized to pyruvate, and lactate is then produced from the pyruvate faster than the body can process it, causing lactate concentrations to rise. The production of lactate is beneficial for NAD+ regeneration (pyruvate is reduced to lactate while NADH is oxidized to NAD+), which is used up in oxidation of glyceraldehyde 3-phosphate during production of pyruvate from glucose, and this ensures that energy production is maintained and exercise can continue. During intense exercise, the respiratory chain cannot keep up with the amount of hydrogen ions that join to form NADH, and cannot regenerate NAD+ quickly enough, so pyruvate is converted to lactate to allow energy production by glycolysis to continue.
The resulting lactate can be used in two ways:
Oxidation back to pyruvate by well-oxygenated muscle cells, heart cells, and brain cells
Pyruvate is then directly used to fuel the Krebs cycle
Conversion to glucose via gluconeogenesis in the liver and release back into circulation by means of the Cori cycle
If blood glucose concentrations are high, the glucose can be used to build up the liver's glycogen stores.
Lactate is continually formed at rest and during all exercise intensities. Lactate serves as a metabolic fuel being produced and oxidatively disposed in resting and exercising muscle and other tissues. Some sources of excess lactate production are metabolism in red blood cells, which lack mitochondria that perform aerobic respiration, and limitations in the rates of enzyme activity in muscle fibers during intense exertion. Lactic acidosis is a physiological condition characterized by accumulation of lactate (especially -lactate), with formation of an excessively high proton concentration [H+] and correspondingly low pH in the tissues, a form of metabolic acidosis.
The first stage in metabolizing glucose is glycolysis, the conversion of glucose to pyruvate− and H+:
When sufficient oxygen is present for aerobic respiration, the pyruvate is oxidized to and water by the Krebs cycle, in which oxidative phosphorylation generates ATP for use in powering the cell.
When insufficient oxygen is present, or when there is insufficient capacity for pyruvate oxidation to keep up with rapid pyruvate production during intense exertion, the pyruvate is converted to lactate− by lactate dehydrogenase), a process that absorbs these protons:
The combined effect is:
The production of lactate from glucose (), when viewed in isolation, releases two H+. The H+ are absorbed in the production of ATP, but H+ is subsequently released during hydrolysis of ATP:
Once the production and use of ATP is included, the overall reaction is
The resulting increase in acidity persists until the excess lactate and protons are converted back to pyruvate, and then to glucose for later use, or to and water for the production of ATP.
Neural tissue energy source
Although glucose is usually assumed to be the main energy source for living tissues, there is evidence that lactate, in preference to glucose, is preferentially metabolized by neurons in the brains of several mammalian species that include mice, rats, and humans. According to the lactate-shuttle hypothesis, glial cells are responsible for transforming glucose into lactate, and for providing lactate to the neurons. Because of this local metabolic activity of glial cells, the extracellular fluid immediately surrounding neurons strongly differs in composition from the blood or cerebrospinal fluid, being much richer with lactate, as was found in microdialysis studies.
Brain development metabolism
Some evidence suggests that lactate is important at early stages of development for brain metabolism in prenatal and early postnatal subjects, with lactate at these stages having higher concentrations in body liquids, and being utilized by the brain preferentially over glucose. It was also hypothesized that lactate may exert a strong action over GABAergic networks in the developing brain, making them more inhibitory than it was previously assumed, acting either through better support of metabolites, or alterations in base intracellular pH levels, or both.
Studies of brain slices of mice show that β-hydroxybutyrate, lactate, and pyruvate act as oxidative energy substrates, causing an increase in the NAD(P)H oxidation phase, that glucose was insufficient as an energy carrier during intense synaptic activity and, finally, that lactate can be an efficient energy substrate capable of sustaining and enhancing brain aerobic energy metabolism in vitro. The study "provides novel data on biphasic NAD(P)H fluorescence transients, an important physiological response to neural activation that has been reproduced in many studies and that is believed to originate predominantly from activity-induced concentration changes to the cellular NADH pools."
Lactate can also serve as an important source of energy for other organs, including the heart and liver. During physical activity, up to 60% of the heart muscle's energy turnover rate derives from lactate oxidation.
Blood testing
Blood tests for lactate are performed to determine the status of the acid base homeostasis in the body. Blood sampling for this purpose is often arterial (even if it is more difficult than venipuncture), because lactate levels differ substantially between arterial and venous, and the arterial level is more representative for this purpose.
During childbirth, lactate levels in the fetus can be quantified by fetal scalp blood testing.
Uses
Polymer precursor
Two molecules of lactic acid can be dehydrated to the lactone lactide. In the presence of catalysts lactide polymerize to either atactic or syndiotactic polylactide (PLA), which are biodegradable polyesters. PLA is an example of a plastic that is not derived from petrochemicals.
Pharmaceutical and cosmetic applications
Lactic acid is also employed in pharmaceutical technology to produce water-soluble lactates from otherwise-insoluble active ingredients. It finds further use in topical preparations and cosmetics to adjust acidity and for its disinfectant and keratolytic properties.
Lactic acid containing bacteria have shown promise in reducing oxaluria with its descaling properties on calcium compounds.
Foods
Fermented food
Lactic acid is found primarily in sour milk products, such as kumis, laban, yogurt, kefir, and some cottage cheeses. The casein in fermented milk is coagulated (curdled) by lactic acid. Lactic acid is also responsible for the sour flavor of sourdough bread.
In lists of nutritional information lactic acid might be included under the term "carbohydrate" (or "carbohydrate by difference") because this often includes everything other than water, protein, fat, ash, and ethanol. If this is the case then the calculated food energy may use the standard per gram that is often used for all carbohydrates. But in some cases lactic acid is ignored in the calculation. The energy density of lactic acid is per 100 g.
Some beers (sour beer) purposely contain lactic acid, one such type being Belgian lambics. Most commonly, this is produced naturally by various strains of bacteria. These bacteria ferment sugars into acids, unlike the yeast that ferment sugar into ethanol. After cooling the wort, yeast and bacteria are allowed to "fall" into the open fermenters. Brewers of more common beer styles would ensure that no such bacteria are allowed to enter the fermenter. Other sour styles of beer include Berliner weisse, Flanders red and American wild ale.
In winemaking, a bacterial process, natural or controlled, is often used to convert the naturally present malic acid to lactic acid, to reduce the sharpness and for other flavor-related reasons. This malolactic fermentation is undertaken by lactic acid bacteria.
While not normally found in significant quantities in fruit, lactic acid is the primary organic acid in akebia fruit, making up 2.12% of the juice.
Separately added
As a food additive it is approved for use in the EU, United States and Australia and New Zealand; it is listed by its INS number 270 or as E number E270. Lactic acid is used as a food preservative, curing agent, and flavoring agent. It is an ingredient in processed foods and is used as a decontaminant during meat processing. Lactic acid is produced commercially by fermentation of carbohydrates such as glucose, sucrose, or lactose, or by chemical synthesis. Carbohydrate sources include corn, beets, and cane sugar.
Forgery
Lactic acid has historically been used to assist with the erasure of inks from official papers to be modified during forgery.
Cleaning products
Lactic acid is used in some liquid cleaners as a descaling agent for removing hard water deposits such as calcium carbonate.
| Physical sciences | Carbon–oxygen bond | null |
172592 | https://en.wikipedia.org/wiki/Organic%20reaction | Organic reaction | Organic reactions are chemical reactions involving organic compounds. The basic organic chemistry reaction types are addition reactions, elimination reactions, substitution reactions, pericyclic reactions, rearrangement reactions, photochemical reactions and redox reactions. In organic synthesis, organic reactions are used in the construction of new organic molecules. The production of many man-made chemicals such as drugs, plastics, food additives, fabrics depend on organic reactions.
The oldest organic reactions are combustion of organic fuels and saponification of fats to make soap. Modern organic chemistry starts with the Wöhler synthesis in 1828. In the history of the Nobel Prize in Chemistry awards have been given for the invention of specific organic reactions such as the Grignard reaction in 1912, the Diels–Alder reaction in 1950, the Wittig reaction in 1979 and olefin metathesis in 2005.
Classifications
Organic chemistry has a strong tradition of naming a specific reaction to its inventor or inventors and a long list of so-called named reactions exists, conservatively estimated at 1000. A very old named reaction is the Claisen rearrangement (1912) and a recent named reaction is the Bingel reaction (1993). When the named reaction is difficult to pronounce or very long as in the Corey–House–Posner–Whitesides reaction it helps to use the abbreviation as in the CBS reduction. The number of reactions hinting at the actual process taking place is much smaller, for example the ene reaction or aldol reaction.
Another approach to organic reactions is by type of organic reagent, many of them inorganic, required in a specific transformation. The major types are oxidizing agents such as osmium tetroxide, reducing agents such as lithium aluminium hydride, bases such as lithium diisopropylamide and acids such as sulfuric acid.
Finally, reactions are also classified by mechanistic class. Commonly these classes are (1) polar, (2) radical, and (3) pericyclic. Polar reactions are characterized by the movement of electron pairs from a well-defined source (a nucleophilic bond or lone pair) to a well-defined sink (an electrophilic center with a low-lying antibonding orbital). Participating atoms undergo changes in charge, both in the formal sense as well as in terms of the actual electron density. The vast majority of organic reactions fall under this category. Radical reactions are characterized by species with unpaired electrons (radicals) and the movement of single electrons. Radical reactions are further divided into chain and nonchain processes. Finally, pericyclic reactions involve the redistribution of chemical bonds along a cyclic transition state. Although electron pairs are formally involved, they move around in a cycle without a true source or sink. These reactions require the continuous overlap of participating orbitals and are governed by orbital symmetry considerations. Of course, some chemical processes may involve steps from two (or even all three) of these categories, so this classification scheme is not necessarily straightforward or clear in all cases. Beyond these classes, transition-metal mediated reactions are often considered to form a fourth category of reactions, although this category encompasses a broad range of elementary organometallic processes, many of which have little in common and very specific.
Fundamentals
Factors governing organic reactions are essentially the same as that of any chemical reaction. Factors specific to organic reactions are those that determine the stability of reactants and products such as conjugation, hyperconjugation and aromaticity and the presence and stability of reactive intermediates such as free radicals, carbocations and carbanions.
An organic compound may consist of many isomers. Selectivity in terms of regioselectivity, diastereoselectivity and enantioselectivity is therefore an important criterion for many organic reactions. The stereochemistry of pericyclic reactions is governed by the Woodward–Hoffmann rules and that of many elimination reactions by Zaitsev's rule.
Organic reactions are important in the production of pharmaceuticals. In a 2006 review, it was estimated that 20% of chemical conversions involved alkylations on nitrogen and oxygen atoms, another 20% involved placement and removal of protective groups, 11% involved formation of new carbon–carbon bond and 10% involved functional group interconversions.
By mechanism
There is no limit to the number of possible organic reactions and mechanisms. However, certain general patterns are observed that can be used to describe many common or useful reactions. Each reaction has a stepwise reaction mechanism that explains how it happens, although this detailed description of steps is not always clear from a list of reactants alone. Organic reactions can be organized into several basic types. Some reactions fit into more than one category. For example, some substitution reactions follow an addition-elimination pathway. This overview isn't intended to include every single organic reaction. Rather, it is intended to cover the basic reactions.
In condensation reactions a small molecule, usually water, is split off when two reactants combine in a chemical reaction. The opposite reaction, when water is consumed in a reaction, is called hydrolysis. Many polymerization reactions are derived from organic reactions. They are divided into addition polymerizations and step-growth polymerizations.
In general the stepwise progression of reaction mechanisms can be represented using arrow pushing techniques in which curved arrows are used to track the movement of electrons as starting materials transition to intermediates and products.
By functional groups
Organic reactions can be categorized based on the type of functional group involved in the reaction as a reactant and the functional group that is formed as a result of this reaction. For example, in the Fries rearrangement the reactant is an ester and the reaction product an alcohol.
An overview of functional groups with their preparation and reactivity is presented below:
Other classification
In heterocyclic chemistry, organic reactions are classified by the type of heterocycle formed with respect to ring-size and type of heteroatom. See for instance the chemistry of indoles. Reactions are also categorized by the change in the carbon framework. Examples are ring expansion and ring contraction, homologation reactions, polymerization reactions, insertion reactions, ring-opening reactions and ring-closing reactions.
Organic reactions can also be classified by the type of bond to carbon with respect to the element involved. More reactions are found in organosilicon chemistry, organosulfur chemistry, organophosphorus chemistry and organofluorine chemistry. With the introduction of carbon-metal bonds the field crosses over to organometallic chemistry.
| Physical sciences | Organic reactions | Chemistry |
172599 | https://en.wikipedia.org/wiki/Search%20and%20rescue | Search and rescue | Search and rescue (SAR) is the search for and provision of aid to people who are in distress or imminent danger. The general field of search and rescue includes many specialty sub-fields, typically determined by the type of terrain the search is conducted over. These include mountain rescue; ground search and rescue, including the use of search and rescue dogs (such as K9 units); urban search and rescue in cities; combat search and rescue on the battlefield and air-sea rescue over water.
International Search and Rescue Advisory Group (INSARAG) is a UN organisation that promotes the exchange of information between national urban search and rescue organisations. The duty to render assistance is covered by Article 98 of the UNCLOS.
Definitions
There are many different definitions of search and rescue, depending on the agency involved and country in question.
Canadian Armed Forces and Canadian Coast Guard: "Search and Rescue comprises the search for, and provision of aid to, persons, ships or other craft which are, or are feared to be, in distress or imminent danger."
United States Coast Guard: "The use of available resources to assist persons or property in potential or actual distress."
United States Department of Defense: A search is "an operation normally coordinated by a Rescue Coordination Center (RCC) or rescue sub-center, using available personnel and facilities to locate persons in distress" and rescue is "an operation to retrieve persons in distress, provide for their initial medical or other needs, and deliver them to a place of safety".
History
One of the world's earliest well-documented SAR efforts ensued following the 1656 wreck of the Dutch merchant ship Vergulde Draeck off the west coast of Australia. Survivors sought help, and in response three separate SAR missions were conducted, without success.
On 29 November 1945, a Sikorsky R-5 performed the first civilian helicopter rescue operation in history, with Sikorsky's chief pilot Dmitry "Jimmy" Viner in the cockpit, using an experimental hoist developed jointly by Sikorsky and Breeze. All five crew members of an oil barge, which had run aground on Penfield Reef, were saved before the barge sank.
In 1983, Korean Air Lines Flight 007 with 269 occupants was shot down by a Soviet aircraft near Sakhalin. The Soviets sent SAR helicopters and boats to Soviet waters, while a search and rescue operation was initiated by U.S., South Korean, and Japanese ships and aircraft in international waters, but no survivors were found.
In July 2009, Air France Flight 447 was lost in the middle of the Atlantic Ocean. An international SAR effort was launched, to no avail. A third effort nearly two years later discovered the crash site and recovered the flight recorders.
In early 2014, Malaysia Airlines Flight 370 crashed under mysterious circumstances. Many nations contributed to the initial SAR effort, which was fruitless. In June 2014, the Australian Transport Safety Bureau commissioned the MV Fugro Equator to lead a three-month survey of the ocean bed. , the search for Flight 370 had become the largest SAR to date.
Types of search and rescue
Ground (lowland) search and rescue
Ground search and rescue is the search for persons who are lost or in distress on land or inland waterways. People may go missing for a variety of reasons. Some may disappear voluntarily, due to issues like domestic abuse. Others disappear for involuntary reasons such as mental illness, getting lost, an accident, death in a location where they cannot be found or, less commonly, due to abduction. Ground search and rescue missions that occur in urban areas should not be confused with "urban search and rescue", which in many jurisdictions refers to the location and extraction of people from collapsed buildings or other entrapments.
In some countries, the police are the primary agency for carrying out searches for a missing person on land. Some places have voluntary search and rescue teams that can be called out to assist these searches.
Search and rescue agencies may contain small specialist teams for executing operations where there are specific environmental risks. Examples include swift water rescue, flood response, technical rope rescue, confined space rescue, over-snow rescue, and thin ice rescue.
Mountain rescue
Mountain rescue relates to search and rescue operations specifically in rugged and mountainous terrain.
Cave rescue
Cave rescue is a highly specialised form of rescue for rescuing injured, trapped or lost cave explorers.
Urban search and rescue
Urban search and rescue (US&R or USAR), also referred to as Heavy Urban Search and Rescue (HUSAR), is the location and rescue of persons from collapsed buildings or other urban and industrial entrapments. Due to the specialised nature of the work, most teams are multi-disciplinary and include personnel from police, fire and emergency medical services. Unlike traditional ground search and rescue workers, most US&R responders also have basic training in structural collapse and the dangers associated with live electrical wires, broken natural gas lines and other hazards. While earthquakes have traditionally been the cause of US&R operations, terrorist attacks and extreme weather such as tornadoes and hurricanes have also resulted in the deployment of these resources.
Combat search and rescue
Combat search and rescue (CSAR) is search and rescue operations that are carried out during war that are within or near combat zones.
Maritime search and rescue
Maritime search and rescue is carried out at sea to save sailors and passengers in distress, or the survivors of downed aircraft. The type of agency which carries out maritime search and rescue varies by country; it may variously be the coast guard, navy or voluntary organisations. When a distressed or missing vessel is located, these organisations deploy helicopters, rescue vessels or any other appropriate vessel to return them to land. In some cases, the agencies may carry out an air-sea rescue (ASR). This refers to the combined use of aircraft (such as flying boats, floatplanes, amphibious helicopters and non-amphibious helicopters equipped with hoists) and surface vessels. Another type of Maritime search and rescue is Submarine rescue. The International Convention on Maritime Search and Rescue (SAR Convention) is the legal framework that applies to international maritime and air-sea rescue.
By nation
Australia
National
The Australian search and rescue service is provided by three authorities; the Joint Rescue Coordination Centre (JRCC) at the Australian Maritime Safety Authority (AMSA), the Australian Defence Force (ADF) and the State/Territory Police Jurisdictions. In a very broad sense, the JRCC respond to national and international registered aircraft, off shore marine incidents and beacon activations. The ADF are responsible for Australian and foreign military personnel, vehicles, vessels and aircraft while within the Australian SRR. Police are responsible for coastal marine incidents, lost persons, unregistered aircraft, inland waterways, ports and identified beacons. The JRCC operates a 24-hour Rescue Coordination Centre (RCC) in Canberra and is responsible for the national coordination of both maritime and aviation search and rescue. The JRCC is also responsible for the management and operation of the Australian ground segment of the Cospas-Sarsat distress beacon detection system. The JRCC's jurisdiction spans Australia and as well as covering 52.8 million square kilometres of the Indian, Pacific and Southern Oceans constituting about 11% of the Earth's surface.
The JRCC is staffed by SAR specialists who have a naval, merchant marine, air force, civil aviation or police service background. The JRCC also coordinates medical evacuations, broadcasts maritime safety information and operates the Australian Ship Reporting System (AUSREP). In coordinating search and rescue missions, the JRCC will call on assistance from organisations as appropriate, such as the Defence forces, Border Protection Command, trained aviation organisations (Civil SAR Units), emergency medical helicopters, state Police services and trained Air Observers from the State Emergency Service. There are also other organisations, such as the non-profit Westpac Life Saver Rescue Helicopter Service that is based at a number of sites around Australia and contracted by various authorities to deliver search and rescue services.
State
State Police in many states operate state-based search and rescue squads, such as the Victoria Police Search and Rescue Squad, which provides specialist expertise, advice and practical assistance in land search and rescue on most terrain including snow and vertical cliff search and rescue. There are also state-based volunteer search and rescue groups such as the NSW SES Bush Search and Rescue in New South Wales and Bush Search and Rescue Victoria in Victoria. These state-based groups draw searchers from bushwalking, mountaineering and specialist rescue clubs within their State. A few groups respond on horseback as mounted search and rescue. The State Emergency Service is a collection of volunteer-based emergency organisations established in each state or territory which are responsible for many rescue efforts in urban and rural areas and in any rescue that results from flood or storm activity. In rural areas the SES conducts most bush search, vertical and road traffic rescues. In urban areas they assist the police and fire services with USAR.
Azerbaijan
Search and rescue operations in Azerbaijan are managed by the Ministry of Emergency Situations onshore in cooperation with the State Civil Aviation Administration in air and the State Maritime Administration offshore.
Belgium
Search and rescue duties along the Belgian part of the North Sea are executed by the Belgian Air Component. From its Koksijde Air Base it operates NH-90 helicopters.
Brazil
Search and rescue duties in Brazil are the responsibility of the Salvarmar Brasil (MRCC Brazil), of the Brazilian Navy and Divisão de Busca e Salvamento (D-SAR) (English: Search and Rescue Division), of the Brazilian Air Force.
Canada
Air and marine Search and rescue duties in Canada are the responsibility of the Canadian Forces and Canadian Coast Guard in conjunction with volunteer organisations. The Department of National Defence (DND) has overall responsibility for the coordinated search and rescue system. SAR operations are organised by Joint Rescue Coordination Centres (JRCC). The JRCC are staffed 24 hours a day by SAR Co-ordinators from the Canadian Coast Guard and Canadian Forces. Authority for the provision of maritime SAR is assigned to the Minister of Fisheries and Oceans by the Canada Shipping Act and the Canada Oceans Act. Ground and inland water search and rescue (GSAR) is the responsibility of provinces and territories with the Royal Canadian Mounted Police (RCMP) and other police forces coordinating operations, often using volunteer GSAR teams operating in specific areas under provincial coordinating bodies.
The Canada Shipping Act, most recently passed in 2001, is the framework document that funds international SAR activities.
The Canadian Forces have five assigned SAR squadrons:
103 Search and Rescue Squadron, CFB Gander, CH-149 Cormorant
413 Transport and Rescue Squadron, CFB Greenwood, CH-149 Cormorant & CC-130 Hercules
424 Transport and Rescue Squadron, CFB Trenton, CH-146 Griffon & CC-130 Hercules
435 Transport and Rescue Squadron, CFB Winnipeg, CC-130 Hercules
442 Transport and Rescue Squadron, CFB Comox, CH-149 Cormorant & CC-115 Buffalo
Plus three Combat Support Squadrons with SAR roles:
417 Combat Support Squadron, CFB Cold Lake, CH-146 Griffon
439 Combat Support Squadron, CFB Bagotville, CH-146 Griffon
444 Combat Support Squadron, CFB Goose Bay, CH-146 Griffon
Some municipalities and provinces have their own SAR units:
Halton Regional Police Service Marine Unit - using marine craft on Lake Ontario
Toronto Police Service Marine Unit - using marine craft on Lake Ontario
Peel Regional Police Marine Unit - using marine craft on Lake Ontario and rivers in Peel Region
Ontario Provincial Police Marine Unit - using marine craft on Great Lakes (excluding Lake Michigan) and Georgian Bay
Durham Regional Police Marine Unit - using marine craft on Lake Ontario and lakes within Durham Region
York Regional Police Marine Unit - using marine craft on Lake Simcoe
Niagara Regional Police Marine Unit - using marine craft on Niagara River and Lake Ontario
Vancouver Police Department - using marine craft on waterways around the City of Vancouver
Heavy Urban Search and Rescue (Toronto) - using land base equipment
Brockville Police Service Marine Patrol Unit - using a boat on the St. Lawrence River
There are also volunteer non-profit associations that conduct SAR in Canada:
British Columbia, there are 80 community based volunteer Groups in B.C. providing GSAR services within assigned areas in conjunction with Police, ambulance and other agencies. The GSAR Groups are represented by the British Columbia Search and Rescue Association
Alberta / BC Cave Rescue, Alberta/British Columbia
Canada Task Force 2, Alberta
Civil Air Search and Rescue Association
ERT Search and Rescue
Grande Prairie Technical Search and Rescue Association, Alberta
Halifax Regional Search and Rescue - Nova Scotia
North Shore Rescue, British Columbia.
Pincher Creek Search and Rescue, Alberta
Québec Secours, Québec.
River Valley Ground Search and Rescue, New Brunswick
Roberts Bank Lifeboat - Delta, BC
Royal Canadian Marine Search and Rescue (RCM SAR)
Sauvetage Bénévole Outaouais - Ottawa Volunteer Search and Rescue - Ottawa, ON and Gatineau, QC
Search and Rescue Manitoba (SARMAN), Manitoba
Vancouver Urban Search and Rescue (Canadian Task Force One), British Columbia
York Sunbury Search & Rescue - New Brunswick
Croatia
In Croatia the SAR Service is part of the Croatian Navy and the Croatian Coast Guard with their headquarter in Rijeka.
Cyprus
The Cyprus Republic Search and Rescue (SAR) system is organised by the Cyprus Joint Rescue Coordination Center (JRCC Larnaca).
The JRCC (Greek: Κέντρο Συντονισμού Έρευνας και Διάσωσης) is an independent agency of the Ministry of Defence of the Republic of Cyprus that started its operations on a 24h basis on 7 August 1995 as a unit of the Cyprus Air Force Command.
On 1 March 2002, the JRCC took full responsibility for investigating, organising, coordinating and executing every SAR incident-operation in the Republic of Cyprus Search & Rescue Region (SRR). JRCC Larnaca operated as a military unit until 26 July 2010, when JRCC was transformed to an independent agency under the Ministry of Defence with the Minister being responsible for its operational aspects. Logistic and technical support is the responsibility of the Ministry of Communications & Works. Its primary mission is to organise the Cyprus Republic Search And Rescue system, to co-ordinate, control and direct SAR operations in its area of responsibility (which is identical to the Nicosia FIR), in order to find and rescue people whose lives are at risk, as a result of an air or naval accident, in the least possible time. This is achieved by coordinating all the different agencies involved such as the Cyprus Police Aviation Unit, the Cyprus Port and Naval Police, the Cyprus National Guard Naval Command, the Cyprus National Guard Air Force Command, the Cyprus Civil Defence and other secondary units.
The JRCC reports directly to the operational control of the Ministry of Defence and it is staffed by qualified personnel of the Cyprus National Guard, mainly from the branches of the Navy and the Air Force.
Northern Cyprus
There are also search and rescue teams in Northern Cyprus. Search and rescue operators in Turkish Republic of North Cyprus are primarily:
SSTB Civil Defence Organisation Presidency (Turkish : Sivil Savunma Teşkilat Başkanlığı)
Emergency Management Committee
DAK Search and Rescue in Natural Disasters (Turkish : Doğal Afetlerde Arama Kurtarma)
AKUT TRNC
Military
TRNC Coast Guard Command
TRNC Coastal Safety (police)
Security Forces Command Search and rescue teams.
Denmark
Search and rescue operators in Denmark are primarily: Danish air force Squadron 722, Danish navy air squadron, naval home guard and the Danish Maritime Safety Administration, coordinated by the Joint Rescue Coordination Centre, operated by the navy and air force in the Danish Naval Commands facilities near Aarhus. Internationally the Danish works mainly with Germany, Norway and Sweden. With the two latter, the annual exercises Baltic SAREX and
Scan-SAR are conducted.
SAR services in Denmark started in 1957 with seven Sikorsky S-55s. Their piston engines produced only and they had limited fuel capacity, so their operational range was short. To increase the operational area, Pembroke twin-engined fixed-wing aircraft were employed for search. These aircraft would localise the distressed person(s) and the S-55s would then rescue them. The SAR service was started for respond to fighter-plane crashes as 79 aircraft crashed, with 62 dead, in the period 1950–1955., but civilian SAR duties are also conducted.
In 1962, eight ship-based Aérospatiale Alouette IIIs were received. These were primarily meant for the ships patrolling the North Atlantic, but also supported the S-55s. In 19641965 the seven S-55s were replaced with eight Sikorsky S-61A helicopters.
In 2007, the Danish Defence held a public display in Horsens, to raise awareness about rescue services and maritime safety. Maritime SAR is important because Denmark has a relative long coast line to its land mass.
In 2008, the SAR forces in Denmark were equipped with eight EH-101, one or two Lynx, 34 naval home guard vessels and 21 rescue vessels, as well as the naval vessels at sea. The EH-101s operate from bases in Aalborg, Skrydstrup and Roskilde. When the sea water temperatures are low a helicopter is also deployed to the island of Bornholm in the Baltic Sea. The Lynx operates from Karup. Maritime vessels are spread out through the entire coastline and on islands. The S-61s and EH-101s have a crew of six: two pilots, a navigator, a flight engineer, a physician and a rescue swimmer.
Estonia
The Estonian Border Guard (Piirivalve) is the Estonian security authority responsible for the border security. It is the main support organisation for search and rescue missions in Estonia, and operates a small fleet of SAR vessels and helicopters.
Finland
In Finland local rescue services (i.e. fire departments) are responsible for land and inland water SAR, the Border Guard is responsible for maritime areas. These organisations alert and decide on the most suitable response for the location and situation. The country also has several volunteer organisations such as the volunteer fire department (VPK), the Finnish Lifeboat Institution (SMPS) and the Red Cross Finland (SPR).
France
The Société Nationale de Sauvetage en Mer (SNSM) provides sea rescue on the French coast and at seas. In 2016, they helped 7,500 people in 5,200 rescues. The service has 41 all-weather rescue boats, 34 first-class rescue boats and 76 second-class lifeboats.
In France, Search and rescue operations are led by different entities according to the rescue area. For sea rescue, the French navy use airborne unit (e.g. Flottille 33F in Brittany) and specialized boats (e.g. "L'abeille Bourbon"). In Mountains, French gendarmerie is equipped with EC-145 'chouka'. In the other areas, French civil protection agency "Securité Civile" works with paramedics, fire unit and hospital mobile unit using EC-145 'Dragon'.
Germany
Search and Rescue in German waters is conducted by the German Maritime Search and Rescue Service with air support by the German Navy, the Federal Police and the German Army Aviation. All incoming requests are coordinated by the Maritime Rescue Coordination Center in Bremen. The DGzRS is a non-governmental organisation entirely supported by donations.
Besides the offshore Search And Rescue services, the German Army Aviation provides 3 SAR Command Posts on a 24/7 basis at Holzdorf Air Base with the Airbus H145 LUH SAR (Light Utility Helikopter Search and Rescue) and at Nörvenich Air Base and 2 at Niederstetten Army Airfield.
Further, the Technisches Hilfswerk is a key component of the German disaster relief framework. It is, among other things, regularly involved in urban search and rescue efforts abroad.
Hong Kong
SAR operations are conducted by the Government Flying Service (GFS) and before 1993 by the Royal Hong Kong Auxiliary Air Force. The GFS conducts maritime SAR within the radius of the Hong Kong Flight Information Region (FIR).
As of 2020, the GFS fleet consists of nine aircraft including:
2 Bombardier Challenger 605 - for aerial SAR surveillance
7 Airbus Helicopters H175 - for inshore and offshore SAR
Other civilian search and rescue units in Hong Kong include:
Civil Aid Service - works in conjunction with the Hong Kong Fire Services Dept and the air support from the Government Flying Service, also provides mountain rescue service
Hong Kong Fire Services/Hong Kong Marine Police - various vessels and rescue divers - with air support from the GFS
Hong Kong Maritime Rescue Co-ordination Centre is responsible for coordinating other civil agencies in regards to marine SAR operations in waters around Hong Kong
Countryside Volunteer Search Team
Iceland
The Icelandic Coast Guard is responsible for coordinating all maritime and aviation search and rescue activities in the Icelandic Search and Rescue Region (SRR), that has the sise of 1.9 million square kilometres. The Icelandic Coast Guard operates JRCC ICELAND in combination with the Coast Guard's operation centre, the maritime traffic service and the coastal radio stations. If aircraft crash site is located on land the control of the rescue operations is diverted to the Icelandic Police, which is responsible for SAR operations on land. The Icelandic Coast Guard (JRCC ICELAND) is the Cospas-Sarsat SAR Point of Contact. ISAVIA, which operates the Air Traffic Control in Iceland, is responsible for the aviation alerting services. The Icelandic Coast Guard operates maritime patrol aircraft, SAR helicopters and patrol vessels.
The Icelandic Association for Search and Rescue (Slysavarnafélagið Landsbjörg) (ICESAR) is a volunteer organisation with about 100 rescue teams located all around the island. ICESAR is a great support to SAR operations both on land and sea. All the rescue teams contain groups of specially trained individuals.
A specialised INSARAG External Classification certified rubble rescue squad operates under the Icelandic Association of Search and Rescue. It was the first rescue squad to arrive in Haiti following the earthquake of 2010.
Indonesia
The National Search and Rescue Agency of Indonesia known in Indonesian as Badan Nasional Pencarian dan Pertolongan abbreviated "BASARNAS", is a government agency responsible for conducting search and rescue duties nationally in Indonesia. BASARNAS may also be assisted in conducting SAR in Indonesia by the TNI, Mobile Brigade Corps, and local Fire brigade units.
Ireland
Maritime SAR services are provided by two civilian bodies - the Irish Coast Guard and the RNLI. The Coast Guard has responsibility for the Irish Search and Rescue Region. The Royal National Lifeboat Institution has 43 lifeboat stations including inland stations at Enniskillen and Lough Derg, the coastguard inshore rescue boats, and community rescue boats at fifteen stations: Ballinskelligs, County Kerry; Ballybunion, County Kerry; Ballyheigue, County Kerry; Banna, County Kerry; Bantry, County Cork; Bunmahon, County Waterford; Cahore, County Wexford; Carna, County Galway; Corrib/Mask Lakes, County Galway; Derrynane, County Kerry; Limerick City (River Shannon); Mallow Search and Rescue, County Cork; Schull, County Cork; Tramore, County Waterford; Waterford City River Rescue; Waterford Marine Search and Rescue. There are some 25 other independent rescue services.
Mountain Rescue in Ireland is provided by 12 voluntary teams based in different regions of the country.
The IRCG operate a number of contracted Sikorsky Search and Rescue helicopters from bases in Dublin, Waterford, Shannon and Sligo under the €500 million contract, from 2010, a previous fleet of Sikorsky S-61N helicopters were replaced with five newer Sikorsky S-92 helicopters. One of the new S-92 helicopters is located at each of the four IRCG bases, with one spare replacement aircraft being rotated between bases.
The Irish Coast Guard are launching a tender for a future SAR Aviation Contract, which is one of several tenders for similar services.
The Irish Air Corps provide top cover for search and rescue over land or sea and is available for maritime and mountain rescue if needed. The Irish Naval Service frequently assists the other agencies in search and rescue. Its patrol ships at sea and the communications centre at Haulbowline maintain a 24-hour watch on all distress frequencies.
Civil Defence Ireland also operates a range of land and inland water search and rescue services.
Israel
SAR in Israel is the responsibility of the IDF Home Front Command Search and Rescue (SAR). The unit was established at its current strength in 1984, combining all the specialist units that were involved with SAR until that time.
The SAR unit is a rapid mobilisation force and has an airborne transport and deployment capability for its personnel and equipment. The unit is composed of reserve personnel, with a regular cadre based at the Bahad 16 Unit training facility. With a focus on urban SAR, the unit operates specialised equipment, including a locally developed device for locating persons trapped under rubble by detecting seismic and acoustic emissions given off by the victims. The SAR unit also uses Search and rescue dogs specially trained to locate people buried under debris.
Israeli SAR resources
Israel Defense Forces
Medical Corps (Israel)
Home Front Command
Bahad 16
Oketz Unit
Israel Police
IsraAid
Magen David Adom
ZAKA
Italy
Italian SAR operations are carried out by the Guardia Costiera, backed up by naval aviation and the air force, including Aeronautica Militare Comando 15° Stormo (15th Wing), the Italian Red Cross, and other organisations.
Jordan
Jordan's Civil Defense Urban Search and Rescue team (USAR) has achieved the UN classification as a heavy USAR team. The team's role mainly earthquake rescue.
Kenya
The Kenya Maritime Authority and the Kenya Civil Aviation Authority are responsible for Aeronautical SAR within Kenya's waterways and aerospace respectively.
Macau
Macau's maritime SAR is conducted by two units:
The Macau Marine Department and responsible for maritime SAR within Macau's waterways. The Macau Search and Rescue Coordination Centre is under the Vessel Traffic Control Centre of Macao of the Maritime Administration of Macau.
Malaysia
For land rescue, Malaysia has two primary SAR units: the Special Malaysia Disaster Assistance and Rescue Team (SMART), which reports to the National Security Council, and the Malaysian Fire and Rescue Department's (FRDM) Special Tactical Operation and Rescue Team of Malaysia (STORM) unit. They are sometimes aided by the jungle experts, the aboriginal police unit named Senoi Praaq, Royal Malaysian Police (RMP) VAT 69 Commando, the Malaysian Armed Forces' special operations force, and the Malaysian Civil Defence Force. Both SMART and STORM, as with other FRDM's special operations, often participate in international SAR missions.
For maritime SAR, it is the responsibility of Malaysia Coast Guard and FRDM, along with support from the RMP's Marine Operations Force and the Royal Malaysian Air Force's PASKAU.
Malta
The responsibility for SAR at sea in the Malta Search and Rescue Region falls under the Armed Forces of Malta (AFM). It is carried out by maritime patrol aircraft, helicopters and vessels under the co-ordination, command and control of the Rescue Co-ordination Centre.
The AFM, in close collaboration with the US Coast Guard, also runs a Search and Rescue Training Centre for International Students in Maritime SAR Mission Co-ordination and Planning. To date more than 30 foreign students from 15 countries including Albania, Cameroon, Croatia, Equatorial Guinea and Kenya have attended these courses.
Malta is also in talks with Libya about enhancing SAR cooperation between the two countries.
Netherlands
SAR responsibility in the Netherlands is held by the Netherlands Coastguard, carried out by vessels and aircraft from various organisations among which mostly the Royal Netherlands Sea Rescue Institution, the Dutch Lifeguard Association, the Ministry of Transport and Water Management and the Ministry of Defence (Netherlands).
New Zealand
New Zealand's Search and Rescue Region extends from the South Pole to the southern border of the Honolulu region, including Norfolk, Tonga, Samoa, and Cook Islands.
Smaller searches are controlled by the local police, who call on LandSAR for land-based operations, such as for lost hikers known as tramping in New Zealand, and the Royal New Zealand Coastguard for coastal maritime incidents. Larger maritime search and rescue events, as well as reports of overdue aircraft, fall under the control of the Rescue Coordination Centre New Zealand (RCCNZ), based in Avalon, which coordinates response from local coastguard, helicopter operators, merchant marine, air force and naval resources.
Urban Search and Rescue falls primarily within the domain of the Fire and Emergency New Zealand, particularly the three USAR Taskforce groups based in Palmerston North, Christchurch, and Auckland. These teams draw together numerous specialists and organisations to achieve an integrated multi-agency response.
Among those organisations that act in a support capacity for FENZ are Response Teams (NZRTs). These are regional rescue groups of professional volunteers that train to a minimum industry standard of USAR Category 1R (USAR Responder), which is also standard for FENZ firefighters. Response Teams are registered with the Ministry of Civil Defence and Emergency Management (MCDEM), and assist their local MCDEM Groups and communities in emergencies to supplement full-time emergency services. Their additional capabilities, which vary among different teams, include: high angle rope rescue, storm response, swift water response, medics, welfare, and rural fire support. Many Response Teams were deployed to assist in the rescue and recovery effort of the 2011 Christchurch earthquake.
Other resources:
Westpac Rescue Helicopter (New Zealand) - charitable organisation
New Zealand Land SAR Search Dogs - the official NZ search dogs group providing land search & rescue services under NZ Land SAR, wilderness and avalanche rescue dogs.
Norway
The search and rescue helicopters are operated by the Royal Norwegian Air Force (RNoAF), who fly 12 Westland Sea Kings. The Norwegian Sea Kings are in the process of being replaced. Between 2020 and 2023 the Sea King fleet will have been phased out by 16 AgustaWestland AW101, named Sar Queen for the purpose. There have been issues with the phasing in of Sar Queen, due to landing sites across the country not being sized for the increased downwash.
Philippines
The agencies responsible for Search and Rescue activities in the Philippines are:
Office of Civil Defense
National Disaster Risk Reduction and Management Council
Philippine Air Force
Philippine Army
Philippine Coast Guard
Bureau of Fire Protection
Portugal
Three different agencies are responsible for providing search and rescue in Portugal. The Portuguese Navy is responsible for all sea rescues, the Portuguese Air Force for all the rescues originating within the airspace, including aircraft crashes and the Autoridade Nacional de Protecção Civil (ANPC) for all inland rescues. All of the above coordinate closely with each other providing a comprehensive search and rescue service.
The Portuguese area of responsibility comprises the Lisbon and Santa Maria Flight Information Regions (FIR).
Poland
In Poland most search and rescue operations are undertaken by the airborne units of the Polish Armed Forces. The Navy currently has the largest SAR fleet of helicopters and also operates a number of small vessels for the purpose of rescuing crewmen of stricken ships. There is also, however a semi-governmental organisation known as the 'Morska Służba Poszukiwania i Ratownictwa' (Maritime Search and Rescue Service) which provides the vast majority of seaborne services to vessels in distress; the service is currently (as of 2010) in the process of overhauling and replacing a large portion of its fleet of lifeboats.
Other civilian search and rescue units in Poland include:
Górskie Ochotnicze Pogotowie Ratunkowe, GOPR (Mountain Volunteer Search and Rescue)
Tatrzańskie Ochotnicze Pogotowie Ratunkowe, TOPR (Tatra Mountains Volunteer Search and Rescue)
Wodne Ochotnicze Pogotowie Ratunkowe, WOPR (Water Volunteer Search and Rescue) - operating on inland and coastal waters
South Africa
Search and Rescue services are offered by various government departments, non governmental organisations, commercial/private organisations and voluntary organisations in South Africa. There is no single organisation responsible for urban, wilderness, swift water, aviation or maritime/sea rescue.
Aviation and maritime incidents are the responsibility of the South African Search and Rescue Organisation (SASAR). SASAR is a voluntary organisation that functions under the auspices of the Department of Transport. Its main role is to search for, assist and carry out rescue operations for the survivors of aircraft or vessel accidents. Depending on the nature of the accident, the RCCs (Aeronatautical Rescue Coordination Centre (ARCC) or Maritime Rescue Coordination Centre (MRCC)) coordinate the search and rescue missions. These operations are carried out by other government departments, non governmental organisations, commercial/private organisations and voluntary organisations.
Local resources:
National Sea Rescue Institute
Wilderness Search and Rescue Cape Town
The Mountain Club of South Africa Search and Rescue
Search and Rescue South Africa
Rescue South Africa
K9 Search and Rescue
Metro Search and Rescue
Spain
Search and rescue duties in Spain are the responsibility of the national government, in conjunction with regional and municipal governments. The Sociedad de Salvamento y Seguridad Marítima is the main organisation, and has overall responsibility for the maritime search and rescue, that also coordinates the SAR efforts with other agencies:
Spanish Navy
Spanish Air and Space Force
Servicio de Vigilancia Aduanera
Servicio Marítimo de la Guardia Civil
Sweden
The Swedish Maritime Administration is responsible for maritime SAR in Swedish waters, and operate seven AgustaWestland AW139 SAR helicopters from five bases along the coast of Sweden. Together with the Swedish Sea Rescue Society and the Swedish Coast Guard they are carrying out Sar in Swedish waters. The Coast guard has 31 larger Ships and 3 De Havilland Canada Dash 8 aeroplanes for SAR. The Swedish Sea Rescue Society is an organisation aiming at saving lives and recovering property at sea. The society operates 68 search and rescue stations and some 185 ships crewed by 2100 volunteers, of those more than 300 are on call at any time, and can respond within 15 minutes. In 2011, the volunteers turned out to an emergency 3274 times. The Swedish Sea Rescue Society is involved 70% of the number SAR missions in Swedish waters.
Switzerland
REGA (Schweizerische REttungsflugwacht / Garde Aérienne / Guardia Aerea) is the air rescue service which provides emergency medical assistance in Switzerland, notably in mountains but also in cases of life-threatening emergencies elsewhere. They will also return a citizen to Switzerland from a foreign country if they are in need of urgent medical care. Rega was established on 27 April 1952 by Dr. Rudolf Bucher, who thought that the Swiss rescue organisation needed a specialised air sub-section.
Taiwan
National Airborne Service Corps (NASC; ) is the agency of the Ministry of the Interior of the Republic of China responsible for executing and providing support for search and rescue, disaster relief, emergency medical service, transportation, monitoring, reconnaissance and patrol in Taiwan.
Coast Guard Administration (CGA; ) is charged with maintaining coastal waters and the pelagic zone patrols, smuggling and stowaway crackdowns, maritime rescues, natural resource conservation, and public services. The CGA is considered a civilian law enforcement agency under the administration of the Executive Yuan, though during emergencies it may be incorporated as part of the Republic of China Armed Forces.
Turkey
Search and Rescue operators in Turkey are primarily:
Civil governmental and non-governmental organisations
Disaster and Emergency Management Presidency; also known as AFAD
AKUT Search and Rescue Association
National Medical Rescue Team (UMKE)
GEA Search and Rescue Association is a group of search and rescue, ecology and social campaigners, founded in 1994, made up of volunteer members.
AKDF Search and Rescue Associations Federation was established in nineteen different districts.
AKA
Search and Rescue and Emergency Aid Association (AKAY)
AKUT
Middle East Search and Rescue, Mountaineering and Outdoor Sports Association (ORDOS)
NAK
National Emergency Search and Rescue Association (NESAR)
Military
Gendarmerie Search and Rescue Battalion Command (JAK)
Gendarmerie Underwater Search and Rescue Teams (SAK)
Diving, Safety, Security, Search and Rescue Team
Air Force Search and Rescue
Coast Guard Command (Turkey) Turkish Coast Guard is also the main Search and Rescue Coordination Authority in Turkish SAR Zone
Ukraine
In Ukraine search and rescue is conducted by the State Search and Rescue Aviation Service of the Ministry of Emergencies of Ukraine Ukraviaposhuk.
United Kingdom
In the UK, land-based searches for a missing person are usually coordinated by the local police. There is a network of local volunteer agencies that can be called out to assist these searches, which are part of the Association of Lowland Search And Rescue. Other voluntary agencies exist to provide specialist search and rescue services, such as the Cave Rescue Organisation and Mountain Rescue Committee of Scotland. These organisations are usually called out indirectly by the police. For example, the British Cave Rescue Council advises that if someone goes missing in a cave, callers should contact the local police who will then summon cave rescue. Urban search and rescue units are run by the fire services.
His Majesty's Coastguard are in charge of maritime search and rescue missions. The Coastguard is one of the four emergency services that can be contacted on 999. Their role is to initiate and coordinate the searches. Lifeboats are provided by volunteer agencies, most often by the Royal National Lifeboat Institution. Aircraft for an air-sea rescue were originally provided by the Royal Navy and Royal Air Force. Under the programme UK-SAR, they are now operated under contract by Bristow Helicopters. The Maritime & Coastguard Agency are launching a tender for their second generation UK search and rescue aviation programme (UKSAR2G), which is one of several tenders for similar services.
Examples of local resources include:
Berkshire Lowland Search and Rescue
Cardiff and Vale Rescue Association
Cave Rescue Organisation
Emergency Response Team Search and Rescue
Mercia Inshore Search and Rescue
Scarborough and Ryedale Mountain Rescue Team
Surrey Search and Rescue
Severn Area Rescue Association
Upper Wharfedale Fell Rescue Association
West Mercia Search and Rescue
United States
In the United States there are many organisations with SAR responsibilities at the national, state and local level. Most day-to-day SAR missions in the US are run by the County Sheriffs, except in states like Alaska, where the State Highway Patrol oversees SAR, or in other areas where SAR services are part of fire/rescue, EMS, or a wholly separate, non-profit organization. They in turn, can request help from other departments, as well as state and national resources if they think they need them. A typical Sheriff's Office has a volunteer SAR team that matches the terrain and population of that county. SAR members are typically trained in the Incident Command System (ICS), first aid, and the outdoor skills needed in that terrain and climate. Most of this article is about the federal response to assist large complicated SAR missions.
In January 2008, the United States Department of Homeland Security (DHS) released the National Response Framework which serves as the guiding document for a federal response during a national emergency. Search and Rescue is divided into four primary elements, while assigning a federal agency with the lead role for each of the four elements.
Structural Collapse-USAR: Department of Homeland Security Federal Emergency Management Agency
Waterborne: United States Coast Guard, United States Coast Guard Auxiliary
Inland-wilderness: United States Department of Interior, National Park Service
Aeronautical: United States Air Force via the Air Force Rescue Coordination Center and USAF rescue wings, groups and squadrons in the Air Combat Command, Pacific Air Forces (for Alaska and Hawaii), Air Education and Training Command, Air Force Reserve Command and the Air National Guard; the Civil Air Patrol in its role as the USAF Auxiliary; and the United States Navy and United States Marine Corps, both Active and Reserve (secondary missions for land-based USN maritime patrol and reconnaissance squadrons and land-based and sea-based USN/USMC helicopter squadrons)
SAR standards adopted by agencies having jurisdiction are developed primarily by non-governmental organisations, including ASTM International and National Fire Protection Association. These standards are adopted also by training and certification organisations such as Mountain Rescue Association and National Association for Search and Rescue to develop training that will meet or exceed those standards. Within ASTM International, standards specific to SAR are developed by Technical Committee F32 on Search and Rescue. Formed in 1988, the committee had 85 current members and jurisdiction of 38 approved standards.
Vietnam
Under command of the Central Government:
National Committee of Search and Rescue is responsible for searching, rescuing and disaster relief.
Central Committee of Prevention of Natural Disasters is responsible for analysing information and monitoring disaster relief processes.
Under command of local People's Committee:
Each province and municipality has a Provincial or City Committee of Prevention of Natural Disaster
Under command of the Ministry of Defense:
General Staff: Department of Rescue of Vietnam People's Army is responsible for coordinating all military rescue activities (including ground force rescue activities).
Navy: Office of Rescue of Vietnam People's Navy is responsible for coordinating naval rescue activities.
Air Force: Office of Rescue of Vietnam People's Air Force is responsible for coordinating air force rescue activities.
Coast Guard: Office of Rescue of Vietnam Coast Guard is responsible for coordinating coastal rescue activities.
Border Guard: Office of Rescue of Vietnam Border Defense Force is responsible for coordinating border rescue activities.
Under command of the Ministry of Public Security:
Vietnam Fire and Rescue Police Department is responsible for fire fighting activities.
Under command of the Ministry of Transport:
Department of Maritime Administration: Vietnam Maritime Search and Rescue Coordination Center (VMRCC) is responsible for maritime rescue activities. VMRCC is divided into 4 Rescue Regions:
Vietnam Maritime Search and Rescue Coordination Center of Region I: operate in Tonkin Gulf
Vietnam Maritime Search and Rescue Coordination Center of Region II: operate in North Central sea
Vietnam Maritime Search and Rescue Coordination Center of Region III: operate in Gulf of Thailand and Southern sea
Vietnam Maritime Search and Rescue Coordination Center of Region IV: operate in South Central sea
Corporation of Air Traffic Management: Vietnam Aviation Search and Rescue Coordination Center (VARCC) is responsible for air rescue activities. VARCC is divided into 3 Rescue Regions:
Vietnam Aviation Search and Rescue Coordination Center of Northern Vietnam: operate in Northern region
Vietnam Aviation Search and Rescue Coordination Center of Central Vietnam: operate in Central region
Vietnam Aviation Search and Rescue Coordination Center of Southern Vietnam: operate in Southern region
Vietnam Railway Rescue and Natural Calamity Response Center of Northern Vietnam: operates in Northern region
Vietnam Railway Rescue and Natural Calamity Response Center of Central Vietnam: operates in Central region
Vietnam Railway Rescue and Natural Calamity Response Center of Southern Vietnam: operates in Southern region
Legitimacy of Search and Rescue
The legitimacy of sea rescue refers to the ethical, legal, and moral justification for acceptance of the act of rescuing individuals or groups of people who are in distress at sea. Sea rescue operations are conducted to save lives, prevent accidents, and offer assistance to those in peril on the water.
Legitimacy in this context encompasses various aspects:
1. Legal Legitimacy: This refers to compliance with international, regional, and national laws and regulations governing maritime safety, search and rescue operations, and the treatment of individuals in distress. These legal frameworks establish a clear foundation for the authority and duty to conduct sea rescue operations.
2. Ethical Legitimacy: Ethical legitimacy pertains to the moral principles and values associated with rescuing people in distress at sea. The duty to save lives, the principle of non-refoulement (not returning individuals to places where they might face persecution), and the broader humanitarian imperative underpin ethical legitimacy.
3. Moral Legitimacy: Moral legitimacy extends beyond legal and ethical considerations and involves public perception and societal acceptance. It is related to how society views and supports rescue efforts. The moral legitimacy of sea rescue operations is often influenced by the belief that saving human lives is a fundamental moral duty.
4. Operational Legitimacy: This aspect concerns the practical and operational effectiveness of sea rescue operations. Ensuring that rescues are conducted efficiently, safely, and with the appropriate resources and expertise contributes to their operational legitimacy.
5. Stakeholder Acceptance: The legitimacy of sea rescue operations also depends on the acceptance and cooperation of various stakeholders, including coastal states, non-governmental organisations (NGOs), international organisations, and local communities. The involvement and support of these actors can enhance the overall legitimacy of rescue efforts.
In the context of the Mediterranean migration crisis, discussions about the legitimacy of sea rescue operations often revolve around these dimensions. Challenges and controversies related to legal ambiguities, accusations against NGOs, pushback practices by some states, and public opinion can all affect the perceived legitimacy of these operations. Balancing legal obligations, ethical imperatives, and the practical challenges of sea rescue remains a complex and ongoing issue, but many argue that the fundamental duty to save lives at sea should be the guiding principle that legitimizes these operations.
The Mediterranean Sea has, for many years, been a major route for migration, hosting numerous asylum seekers, refugees, and economic migrants seeking a better life in Europe. This region is known for its perilous journey, often undertaken in overcrowded and unseaworthy vessels, leading to frequent humanitarian crises at sea. As a result, the question of the legal legitimacy of sea rescue operations in the Mediterranean has emerged as a central concern.
2. International Law of the Sea
International law establishes a clear foundation for the legitimacy of sea rescue operations in the Mediterranean.
The United Nations Convention on the Law of the Sea (UNCLOS) and international conventions like the Safety of Life at Sea (SOLAS) and the International Convention on Maritime Search and Rescue (SAR) obligate states to render assistance to those in distress at sea. These conventions provide a solid legal basis for sea rescue operations, emphasizing the duty to save lives.
2.1 United Nations Convention on the Law of the Sea (UNCLOS)
The United Nations Convention on the Law of the Sea (UNCLOS) provides a legal framework for various maritime activities, including sea rescue operations. This framework can be introduced as a two-step process. First, states are sorted into groups depending on activity and typical interest in the sea, principally coastal states and flag states. Second, the sea is divided into different zones (the territorial sea, the contiguous zone, the exclusive economic zone, the continental shelf, and the high seas) where different categories of states have different authority.
UNCLOS does establish the legal basis for conducting such operations by outlining key principles and obligations related to maritime safety and search and rescue at sea. These principles contribute to the legitimacy of sea rescue efforts:
1. Duty to Render Assistance: UNCLOS, under its article 98, places an obligation on all vessels and aircraft to render assistance to any person found at sea in danger of being lost and to inform the appropriate authorities. This duty underscores the moral and legal imperative to provide assistance to those in distress at sea.
2. Non-refoulment principle: While UNCLOS doesn't explicitly mention the principle of non-refoulement, it is a fundamental aspect of international refugee law and human rights law, which plays a significant role in the legitimacy of sea rescue. Non-refoulement prohibits returning individuals to places where their lives or freedoms might be threatened.
3. SAR Regions and Coordination: UNCLOS encourages the establishment of search and rescue (SAR) regions and coordination centres. These provisions promote efficient and effective responses to maritime distress situations, enhancing the legitimacy of rescue efforts.
2.1.1 Duty to Render Assistance
The duty to render assistance at sea is set out by Article 98 of the UNCLOS, as follows:
‘1. Every State shall require the master of a ship flying its flag, in so far as he can do so without serious danger to the ship, the crew or the passengers:
(a) to render assistance to any person found at sea in danger of being lost;
(b) to proceed with all possible speed to the rescue of persons in distress, if informed of their need of assistance, in so far as such action may reasonably be expected of him; (…)’
It imposes a general duty on state parties to require their vessels “to render assistance to any person found at sea in danger of being lost” and “to proceed with all possible speed to the rescue of persons in distress, if informed of their need of assistance, in so far as such action may reasonably be expected”. Where ships collide, each ship’s master is to be obliged “to render assistance to the other ship, its crew and its passengers and, where possible, to inform the other ship of the name of his own ship, its port of registry and the nearest port at which it will call”. No immigration or other exceptions apply.
UNCLOS also at Article 98(2) imposes a duty on state parties to “promote the establishment, operation and maintenance of an adequate and effective search and rescue service regarding safety on and over the sea and, where circumstances so require, by way of mutual regional arrangements cooperate with neighbouring States for this purpose".
On the face of it, UNCLOS Article 98 obliges state parties, including the United Kingdom, to rescue refugees in distress at sea and also to operate a search and rescue system. There is an argument that Article 98 falls within the part of UNCLOS addressed only to the high seas and therefore the obligation does not apply within territorial waters or the contiguous zone (see Article 86). The words of Article 98 itself do not support this approach.
Article 18 seems to imply a duty of rescue in territorial waters and the contiguous zone; and, in any event, other conventions apply similar obligations without any similar restriction.
2.1.2 Non-refoulement Principle
The principle of non-refoulement refers to the obligation on states not to send individuals to territories in which they may be persecuted, or in which they are at risk of torture or other serious harm. It may not immediately correlate with the right of every one to seek asylum, but it does clearly place limits on what states may lawfully do.
This rule is solidly grounded in international human rights and refugee law, in treaty, in doctrine, and in customary international law. It is an inherent aspect of the absolute prohibition of torture, even sharing perhaps in some of the latter's jus cogens character. It applies independently of any formal recognition of refugee status or entitlement to other forms of protection, and it applies to the actions of states, wherever undertaken, whether at the land border, or in maritime zones, including the high seas.
Its essential characteristics are acts attributable to the state or other international actor, which have the foreseeable effect of exposing the individual to a serious risk of irreversible harm, contrary to international law.
UNHCR's Executive Committee, indeed, has particularly emphasized the importance of fully respecting the principle of non-refoulement in the context of maritime operations:
‘… interception measures should not result in asylum-seekers and refugees being denied access to international protection, or in those in need of international protection being returned, directly or indirectly, to the frontiers of territories where their life or freedom would be threatened on account of a Convention ground, or where the person has other grounds for protection based on international law.’
2.2 International Convention on Maritime Search and Rescue (SAR)
UNCLOS has been described as a “quasi-constitution for the oceans”. Three (or maybe four) further conventions set out the details of search and rescue obligations.
The first of these is the International Convention on Maritime Search and Rescue (SAR), which was adopted by the International Maritime Organisation (IMO) in 1979 and came into force in 1985. The International Convention on Maritime Search and Rescue aims to enhance the effectiveness of search and rescue operations at sea by establishing a framework for coordination, cooperation, and the provision of assistance to persons in distress. It aims to create an international system for coordinating rescue operations that guarantees their efficiency and safety. States parties are thus invited to conclude SAR agreements with neighbouring states to regulate and coordinate SAR operations and services in the agreed maritime zone. Such agreements technically and operationally implement the obligation set out in Article 98 (2) of the UNCLOS, which provides that, where needed, neighbouring states shall cooperate through regional agreements to promote and maintain adequate and effective SAR services.
Following an amendment taking effect in 2004 for all parties except Malta, which formally objected, the obligation to provide assistance to a person at distress at sea expressly applies “regardless of the nationality or status of such a person or the circumstances in which that person is found” (para 2.1.10). Once a person has been rescued, they must be delivered to a “place of safety” (para 1.3.2 and 3.1.9). In particular, the Maritime Safety Committee (MSC) of the International Maritime Organisation (IMO) adopted two resolutions that amended SAR Convention (and SOLAS Convention), and which entered into force 1 July 2006. Consequently, Article 3 (1) (9) of the SAR Convention now provides:
‘Parties shall co-ordinate and co-operate to ensure that masters of ships providing assistance by embarking persons in distress at sea are released from their obligations with minimum further deviation from the ships’ intended voyage (…). The Party responsible for the search and rescue region in which such assistance is rendered shall exercise primary responsibilityfor ensuring such co-ordination and co-operation occurs, so that survivors assisted are disembarked from the assisting ship and delivered to a place of safety (…). In these cases, the relevant Parties shall arrange for such disembarkation to be effective as soon as reasonably practicable’.
According to the MSC Guidelines,a ‘place of safety’ means a location where the rescue operations can be considered as completed. In accordance with Principle 6.14 of the Guidelines, the rescue unit can be the place of safety, but only provisionally. In fact, the text insists on the role that the flag state and the coastal state should play in substituting for the master of the rescuing vessel (Principle 6.13).
Moreover, pursuant to the same guidelines, the state in whose SAR zone the operation took place has the duty to provide or, at least, to secure a place of safety for the rescued persons (Principle 2.5). This Principle simply requires that the coastal state carries out the SAR operations and brings them effectively to an end, i.e., not leaving the rescued persons (whatever their status) at sea. Considering that the MSC Guidelines are not binding, Principle 2.5 suggests that the coastal state has a ‘residual obligation’ to allow disembarkation on its own territory when it has not been possible to do so safely anywhere else.[16] This has been clarified by the IMO Facilitation Committee (FAL), which adopted the ‘Principles relating to administrative procedures for disembarking persons rescued at sea’.
Another issue concerns the identification of the beneficiary of the obligation: is the duty to render assistance a purely inter-state obligation or does it entail a right to be rescued for people in distress at sea? The law of the sea is a field of international law where individuals or more generally private actors have little involvement. The main aim of the law of the sea consists of allocating obligations and rights in different maritime zones to states. However, the multiplication of activities at sea and the increased human presence led to the question of the protection of the human element, in particular of the application of human rights at sea. Many scholars have already discussed the application of the relevant human rights treaties at sea, pointing out that the law of the sea, specifically the LOSC, pursues some community interests, among which the protection of human rights. Building on this scholarship, the duty to render assistance can be considered to be the operational obligation deriving from the application of the human right to life at sea.
2.3 International Convention for the Safety of Life at Sea (SOLAS)
The second one is the International Convention for the Safety of Life at Sea of 1974 (SOLAS), which is primarily concerned with the seaworthiness of ships. This convention is also widely ratified, including by the United Kingdom. It is a very substantial and detailed document, the first version of which was adopted in 1914 in response to the sinking of the Titanic.
SOLAS imposes a duty on the master of a ship at sea, on receiving information from any source that persons are in distress at sea, “to proceed with all speed to their assistance” (Chapter V, regulation 33-1 as amended). While the wording appears directed to the master in person, the obligation is probably better read as being on the state party to ensure that a master of a ship acts in the required way.
The convention also imposes on state parties an obligation “to ensure that necessary arrangements are made for distress communication and co-ordination in their area of responsibility and for the rescue of persons in distress at sea around its coasts” (Chapter V, regulation 7). These duties apply to “all ships on all voyages” other than ships of war and ships in the Great Lakes and tributaries (Chapter V, regulation 1).
Mirroring the 2004 adjustments to SAR, SOLAS was also amended explicitly to state that the duty to provide assistance “applies regardless of the nationality or status of such persons or the circumstances in which they are found”. It further provides that a rescued person should be treated “with humanity” while on board a rescue ship and then delivered to a place of safety.
2.4 Salvage conventions
Finally, there are two international conventions on salvage, both of which impose a requirement to render assistance at sea. The reason the salvage conventions impose this obligation is to ensure that priority is given to saving lives rather than property in a salvage situation.
The first is the Brussel Convention for the Unification of Certain Rules with Respect to Assistance and Salvage at Sea of 1910, or simply the Brussels Convention. Article 11 provides that:
‘Every master is bound, so far as he can do so without serious danger to his vessel, her crew and her passengers, to render assistance to everybody, even though an enemy, found at sea in danger of being lost.
The second is the International Convention on Salvage of 1989 (the Salvage Convention). Article 10 requires every ship’s master, “so far as he can do so without serious danger to his vessel and persons thereon, to render assistance to any person in danger of being lost at sea” and imposes a duty on state parties to adopt the measures necessary to enforce that duty.
3. European Union Law
The European Union (EU) has played a significant role in shaping the legal framework for sea rescue in the Mediterranean.
First of all, it has established the European Border and Coast Guard Agency (Frontex) to coordinate and support member states in border and maritime security.
3.1 Frontex
Frontex is the European Union agency responsible for “European integrated border management” (EU, 2016). Its mission is to “ensure safe and well-functioning external borders providing security” (Frontex, 2019).
Frontex was founded by a Council Regulation in 2004 as the “European Agency for the Management of Operational Cooperation at the External Borders” and became operational a year later. The founding regulation was amended several times, until it was replaced by a new Regulation in 2016, establishing a “European Border and Coast Guard” (European Commission, 2004; EU, 2016). These repeated legal revisions are mirrored by a rapid organisational growth of the agency. Frontex, currently, has a budget of over 300 million euros and will dispose of 11.3 billion euros in 2021–2027, with which it aims to finance a 10,000 standing corps of operational staff and acquire its own ships, vehicles and planes (European Commission, 2018).Frontex plays a central role in managing and regulating sea rescue operations within the EU. EU law, particularly regulations related to Frontex, outlines the responsibilities, actions, and legal parameters for Frontex-coordinated rescue operations.
Ten studies (14 per cent) shed light on the nature and implementation of Frontex’ activities. The Frontex Regulation has been repeatedly revised to expand the agency’s mandate. In fact, Article 8 of the EU Regulation establishing the European Border and Coast Guard identified no less than 21 tasks that Frontex is mandated to carry out. Beyond merely listing these responsibilities, these studies describe how such official tasks are translated into practice.
Frontex can be viewed as the EU’s attempt to implement “integrated border management” (IBM), which aims to introduce non-arrival measures and curb illegal migration. This requires cooperation of Frontex with various partners, such as Member States, third countries and relevant other agencies (e.g. Europol) (Demmelhuber, 2011). Frontex is primarily preoccupied with joint operations to reinforce border control, many of which focus on the Mediterranean (Üstübici and Içduygu, 2018). To this end, Frontex coordinates the deployment of border officials, who participate in border patrol, thereby contributing to the arrest of suspected facilitators and setting in motion a continuous decline in the number of migrants crossing (Tryfon, 2012). An analysis of hot spot “Moria” on Lesbos, Greece, demonstrates that Frontex personnel not only monitors the border but also engages in migrant identification and registration. As migrants arrive in the camp, they access a “registration street”, in which Frontex screeners identify the nationality of migrants, while advanced-level document officers check the veracity of their identification papers. Fingerprinters take fingerprints and enter these into databases, after which Frontex’ partners begin the asylum process and provide humanitarian aid.
To facilitate the effective and efficient deployment of personnel and capabilities, Frontex also gathers intelligence and produces its own risk analyses, which address the (expected) situation at the EU external borders. As its intelligence role is growing, the agency collects data from cross-border movements inside the EU and from entrance information that Member States and the various organisational partners provide (Den Boer, 2015; Takle, 2017). Some of this information is also retrieved from third countries, with which Frontex has concluded working arrangements. This is beneficial for Frontex as it ensures continuous data exchange on relevant migration trends and the sharing of best practices (Coman-Kund, 2018). In turn, representatives of these third countries have been participant observers in Frontex operations (Sagrera, 2014). To what extent these agreements are binding remains subject to debate (cf. Sagrera, 2014; Coman-Kund, 2018).
Finally, Frontex also mediates between the Member States to coordinate joint returns when multiple Member States have nationals of the same country to be sent back. The agency makes sure that returnees are gathered in one Member State from where a collective return flight departs. A Frontex staff member monitors compliance with its Code of Conduct when it organises the return flight (Pirjola, 2015).
It is striking how rapid Frontex’ responsibilities have expanded. Frontex’ tasks, first of all, horizontally expanded with its mandate to set up rapid response teams that can be deployed on short notice. Frontex’ tasks have also deepened as existing responsibilities were elaborated. For instance, Frontex received the authority to organise return flights. At this moment, Frontex could even contribute to the reintroduction of border controls when a Member State does not comply with the recommendations that follow its vulnerability assessments (Scipioni, 2017).
Current studies on Frontex display strong links to the field of critical theory, migration law research and literature on EU public administration. This has provided many insightful publications. Yet, there are very few connections with organisation and management literature, crisis studies and policing scholarship, even though these disciplines may improve our understanding of Frontex’ agency characteristics, how it operates during joint operations and how Frontex negotiates tensions between providing care and exercising control. Likewise, EU studies can shed light on the relations between EU agencies and other European actors and what this means for their relative autonomy (e.g. Egeberg and Trondal, 2011). Next, while there are some links between the five themes, these connections can be reinforced in the interest of theory-building. For instance, there is a suggestion that Frontex’ autonomy (in line with neo-institutionalist explanations) reduces its accountability (Carrera et al., 2013). Likewise, Frontex’ limited accountability and transparency may help to explain its continuing struggles with human rights. In general, as Frontex continues to expand, new theory-building is needed to further assess and understand the agency and its impact.
3.2 Common European Asylum System
The European Union has established a Common European Asylum System (CEAS).
It is a legal and policy framework developed to guarantee harmonised and uniform standards for people seeking international protection in the EU. It is based on an understanding that the EU, an area of open borders and freedom of movement where countries share the same fundamental values, needs to have a common approach to implement transparent, effective and equitable procedures.
CEAS emphasizes a shared responsibility to process applicants for international protection in a dignified manner, ensuring fair treatment and similar procedures in examining cases, irrelevant of the country where the application is lodged. At its core, CEAS aims to achieve: a clear functional process to determine which country is responsible for examining an application for protection; a set of common standards to inform fair and efficient asylum procedures; a set of common minimum conditions for the dignified reception of applicants for protection and convergence on the criteria for granting protection associated with those statuses.
Within the context of CEAS, the Tampere Declaration set out the foundation for a comprehensive approach to migration by addressing political, human rights and developmental issues in countries and regions of origin and transit.68 Through this agreement with governments, legislative and policy measures were adopted at the EU level to set a framework to manage high influxes of displaced persons by accommodating persons in need of protection while supporting Member States experiencing pressure on their asylum systems.
After the first phase (from 1999 to 2005), Member States reflected on the functioning of CEAS and implemented improvements to the five legislations that govern the minimum standards of the European asylum system.
The increased – and often uneven – pressure that national asylum and reception systems in EU+ countries faced since 2015 presented both a challenge and an opportunity for EU+ countries to take bold steps toward systemic and commonly-agreed solutions for further harmonisation, on the basis of solidarity and responsibility-sharing. Above all, it underlined the importance of the very existence of CEAS and a common migration policy – to have an EU-wide framework to manage mixed migratory flows,iv including border management, international protection and the return of rejected applicants. In the EU context, mixed migratory flows are defined as "complex migratory population movements, including refugees, asylum seekers, economic migrants and other types of migrants as opposed to migratory population movements that consist entirely of one category of migrants".
3.3 Charter of Fundamental Rights of the European Union
EU law is bound by the Charter of Fundamental Rights of the European Union, which brings together the most important personal freedoms and rights enjoyed by citizens of the EU into one legally binding document.
The Charter was declared in 2000, and came into force in December 2009 along with the Treaty of Lisbon. It includes provisions related to human dignity, non-discrimination, and asylum. The Convention is divided into 6 chapters (Dignity, Freedoms, Equality, Solidarity, Citizen’s rights, Justice) and articles. These rights and principles are instrumental in shaping the legitimacy of sea rescue operations.
3.4 Schengen Area
Moreover, for EU member states participating in the Schengen Area, there are legal obligations and responsibilities for border control and management, including responding to situations at sea that may require rescue operations.
The border-free Schengen Area guarantees free movement to more than 400 million EU citizens, along with non-EU nationals living in the EU or visiting the EU as tourists, exchange students or for business purposes (anyone legally present in the EU). Free movement of persons enables every EU citizen to travel, work and live in an EU country without special formalities. Schengen underpins this freedom by enabling citizens to move around the Schengen Area without being subject to border checks.
Today, the Schengen Area encompasses most EU countries, except for Bulgaria, Cyprus, Ireland and Romania. However, Bulgaria and Romania are currently in the process of joining the Schengen Area and already applying the Schengen acquis to a large extent. On 1 January 2023 Croatia became the newest member state to join the Schengen area. Additionally, also the non-EU States Iceland, Norway, Switzerland and Liechtenstein have joined the Schengen Area.
Aircraft
Rotary and fixed wing aircraft are used for air and sea rescue. A list of common aircraft used:
Aérospatiale SA330 Puma
Aérospatiale SA360 Dauphin
Airbus Helicopters H175
AgustaWestland AW109
AgustaWestland AW139
AgustaWestland AW101
AgustaWestland CH-149 Cormorant
Bell UH-1 Iroquois
Bell CH-146 Griffon
Boeing Vertol CH-46 Sea Knight
CH-113 Labrador and 113A Voyageur
Eurocopter Dauphin - variant of Aérospatiale SA 360 Dauphin
Eurocopter Dolphin HH-65
Eurocopter AS365 Dauphin 2
Eurocopter AS332 Super Puma
Eurocopter EC225 Super Puma
Eurocopter AS532 Cougar
Lockheed HC-130 Hercules
Lockheed P-3 Orion
Sikorsky S-61
Sikorsky S-70 Blackhawk
Sikorsky HH-60 Jayhawk
Sikorsky HH-60 Pave Hawk
Sikorsky SH-60 Seahawk
Sikorsky S-76
Sikorsky S-92
Sikorsky H-92 Superhawk
Sikorsky CH-148 Cyclone
Westland Sea King
Westland Wessex HC2
| Biology and health sciences | General concepts | Health |
172644 | https://en.wikipedia.org/wiki/Handedness | Handedness | In human biology, handedness is an individual's preferential use of one hand, known as the dominant hand, due to it being stronger, faster or more dextrous. The other hand, comparatively often the weaker, less dextrous or simply less subjectively preferred, is called the non-dominant hand. In a study from 1975 on 7,688 children in US grades 1–6, left handers comprised 9.6% of the sample, with 10.5% of male children and 8.7% of female children being left-handed. Overall, around 90% of people are right-handed. Handedness is often defined by one's writing hand, as it is fairly common for people to prefer to do a particular task with a particular hand. There are people with true ambidexterity (equal preference of either hand), but it is rare—most people prefer using one hand for most purposes. However, in some cultures, the use of the left hand can be considered disrespectful.
Most of the current research suggests that left-handedness has an epigenetic marker—a combination of genetics, biology and the environment.
Because the vast majority of the population is right-handed, many devices are designed for use by right-handed people, making their use by left-handed people more difficult. In many countries, left-handed people are or were required to write with their right hands. However, left-handed people have an advantage in sports that involve aiming at a target in an area of an opponent's control, as their opponents are more accustomed to the right-handed majority. As a result, they are over-represented in baseball, tennis, fencing, cricket, boxing, and mixed martial arts.
Types
Right-handedness is the most common type. Right-handed people are more skillful with their right hands. Studies suggest that approximately 90% of people are right-handed.
Left-handedness is less common. Studies suggest that approximately 10% of people are left-handed.
Ambidexterity refers to having equal ability in both hands. Those who learn it still tend to favor their originally dominant hand. This is uncommon, with about a 1% prevalence.
Mixed-handedness or cross-dominance is the change of hand preference between different tasks. This is about as widespread as left-handedness. This is highly associated with the person's childhood brain development.
Measurement
Handedness may be measured behaviourally (performance measures) or through questionnaires (preference measures). The Edinburgh Handedness Inventory has been used since 1971 but contains some dated questions and is hard to score. Revisions have been published by Veale and by Williams. The longer Waterloo Handedness Questionnaire is not widely accessible. More recently, the Flinders Handedness Survey (FLANDERS) has been developed.
Evolution
Some non-human primates have a preferred hand for tasks, but they do not display a strong right-biased preference like modern humans, with individuals equally split between right-handed and left-handed preferences. When exactly a right handed preference developed in the human lineage is unknown, though it is known through various means that Neanderthals had a right-handedness bias like modern humans. Attempts to determine handedness of early humans by analysing the morphology of lithic artefacts have been found to be unreliable.
Causes
There are several theories of how handedness develops.
Genetic factors
Handedness displays a complex inheritance pattern. For example, if both parents of a child are left-handed, there is a 26% chance of that child being left-handed. A large study of twins from 25,732 families by Medland et al. (2006) indicates that the heritability of handedness is roughly 24%.
Two theoretical single-gene models have been proposed to explain the patterns of inheritance of handedness, by Marian Annett of the University of Leicester, and by Chris McManus of UCL.
However, growing evidence from linkage and genome-wide association studies suggests that genetic variance in handedness cannot be explained by a single genetic locus. From these studies, McManus et al. now conclude that handedness is polygenic and estimate that at least 40 loci contribute to the trait.
Brandler et al. performed a genome-wide association study for a measure of relative hand skill and found that genes involved in the determination of left-right asymmetry in the body play a key role in handedness. Brandler and Paracchini suggest the same mechanisms that determine left-right asymmetry in the body (e.g. nodal signaling and ciliogenesis) also play a role in the development of brain asymmetry
(handedness being a reflection of brain asymmetry for motor function).
In 2019, Wiberg et al. performed a genome-wide association study and found that handedness was significantly associated with four loci, three of them in genes encoding proteins involved in brain development.
Prenatal hormone exposure
Four studies have indicated that individuals who have had in-utero exposure to diethylstilbestrol (a synthetic estrogen-based medication used between 1940 and 1971) were more likely to be left-handed over the clinical control group. Diethylstilbestrol animal studies "suggest that estrogen affects the developing brain, including the part that governs sexual behavior and right and left dominance".
Ultrasound
Another theory is that ultrasound may sometimes affect the brains of unborn children, causing higher rates of left-handedness in children whose mothers receive ultrasound during pregnancy. Research suggests there may be a weak association between ultrasound screening (sonography used to check the healthy development of the fetus and mother) and left-handedness.
Epigenetic markers
Twin studies indicate that genetic factors explain 25% of the variance in handedness, and environmental factors the remaining 75%. While the molecular basis of handedness epigenetics is largely unclear, Ocklenburg et al. (2017) found that asymmetric methylation of CpG sites plays a key role for gene expression asymmetries related to handedness.
Language dominance
One common handedness theory is the brain hemisphere division of labor. In most people, the left side of the brain controls speaking. The theory suggests it is more efficient for the brain to divide major tasks between the hemispheres—thus most people may use the non-speaking (right) hemisphere for perception and gross motor skills. As speech is a very complex motor control task, the specialised fine motor areas controlling speech are most efficiently used to also control fine motor movement in the dominant hand. As the right hand is controlled by the left hemisphere (and the left hand is controlled by the right hemisphere) most people are, therefore right-handed. The theory depends on left-handed people having a reversed organisation. However, the majority of left-handers have been found to have left-hemisphere language dominance—just like right-handers. Only around 30% of left-handers are not left-hemisphere dominant for language. Some of those have reversed brain organisation, where the verbal processing takes place in the right-hemisphere and visuospatial processing is dominant to the left hemisphere. Others have more ambiguous bilateral organisation, where both hemispheres do parts of typically lateralised functions. When tasks designed to investigate lateralisation (preference for handedness) are averaged across a group of left-handers, the overall effect is that left-handers show the same pattern of data as right-handers, but with a reduced asymmetry. This finding is likely due to the small proportion of left-handers who have atypical brain organisation. The majority of the evidence comes from literature assessing oral language production and comprehension. When it comes to writing, findings from recent studies were inconclusive for a difference in lateralization for writing between left-handers and right-handers.
Developmental timeline
Researchers studied fetuses in utero and determined that handedness in the womb was a very accurate predictor of handedness after birth. In a 2013 study, 39% of infants (6 to 14 months) and 97% of toddlers (18 to 24 months) demonstrated a hand preference.
Infants have been observed to fluctuate heavily when choosing a hand to lead in grasping and object manipulation tasks, especially in one- versus two-handed grasping. Between 36 and 48 months, there is a significant decline in variability between handedness in one-handed grasping; it can be seen earlier in two-handed manipulation. Children of 18–36 months showed more hand preference when performing bi-manipulation tasks than with simple grasping.
The decrease in handedness variability in children of 36–48 months may be attributable to preschool or kindergarten attendance due to increased single-hand activities such as writing and coloring. Scharoun and Bryden noted that right-handed preference increases with age up to the teenage years.
Correlation with other factors
The modern turn in handedness research has been towards emphasizing degree rather than direction of handedness as a critical variable.
Intelligence
In his book Right-Hand, Left-Hand, Chris McManus of University College London argues that the proportion of left-handers is increasing, and that an above-average quota of high achievers have been left-handed. He says that left-handers' brains are structured in a way that increases their range of abilities, and that the genes that determine left-handedness also govern development of the brain's language centers.
Writing in Scientific American, he states:
Studies in the U.K., U.S. and Australia have revealed that left-handed people differ from right-handers by only one IQ point, which is not noteworthy ... Left-handers' brains are structured differently from right-handers' in ways that can allow them to process language, spatial relations and emotions in more diverse and potentially creative ways. Also, a slightly larger number of left-handers than right-handers are especially gifted in music and math. A study of musicians in professional orchestras found a significantly greater proportion of talented left-handers, even among those who played instruments that seem designed for right-handers, such as violins. Similarly, studies of adolescents who took tests to assess mathematical giftedness found many more left-handers in the population.
Left-handers are overrepresented among those with lower cognitive skills and mental impairments, with those with intellectual disability being roughly twice as likely to be left-handed, as well as generally lower cognitive and non-cognitive abilities amongst left-handed children. Left-handers are nevertheless also overrepresented in high IQ societies, such as Mensa. A 2005 study found that "approximately 20% of the members of Mensa are lefthanded, double the proportion in most general populations".
Ghayas & Adil (2007) found that left-handers were significantly more likely to perform better on intelligence tests than right-handers and that right-handers also took more time to complete the tests. In a systematic review and meta-analysis, Ntolka & Papadatou-Pastou (2018) found that right-handers had higher IQ scores, but that difference was negligible (about 1.5 points).
The prevalence of difficulties in left-right discrimination was investigated in a cohort of 2,720 adult members of Mensa and Intertel by Storfer. According to the study, 7.2% of the men and 18.8% of the women evaluated their left-right directional sense as poor or below average; moreover participants who were relatively ambidextrous experienced problems more frequently than did those who were more strongly left- or right-handed. The study also revealed an effect of age, with younger participants reporting more problems.
Early childhood intelligence
Nelson, Campbell, and Michel studied infants and whether developing handedness during infancy correlated with language abilities in toddlers. In the article they assessed 38 infants and followed them through to 12 months and then again once they became toddlers from 18 to 24 months. They discovered that when a child developed a consistent use of their right or left hand during infancy (such as using the right hand to put the pacifier back in, or grasping random objects with the left hand), they were more likely to have superior language skills as a toddler. Children who became lateral later than infancy (i.e., when they were toddlers) showed normal development of language and had typical language scores. The researchers used Bayley scales of infant and toddler development to assess the subjects.
Music
In two studies, Diana Deutsch found that left-handers, particularly those with mixed-hand preference, performed significantly better than right-handers in musical memory tasks. There are also handedness differences in perception of musical patterns. Left-handers as a group differ from right-handers, and are more heterogeneous than right-handers, in perception of certain stereo illusions, such as the octave illusion, the scale illusion, and the glissando illusion.
Health
Studies have found a positive correlation between left-handedness and several specific physical and mental disorders and health problems, including:
Lower birth weight and complications at birth are positively correlated with left-handedness.
A variety of neuropsychiatric and developmental disorders like autism spectrum, bipolar disorder, anxiety disorders, schizophrenia, and alcoholism have been associated with left- and mixed-handedness.
A 2012 study showed that nearly 40% of children with cerebral palsy were left-handed, while another study demonstrated that left-handedness was associated with a 62% increased risk of Parkinson's disease in women, but not in men. Another study suggests that the risk of developing multiple sclerosis increases for left-handed women, but the effect is unknown for men at this point.
Left-handed women may have a higher risk of breast cancer than right-handed women and the effect is greater in post-menopausal women.
At least one study maintains that left-handers are more likely to suffer from heart disease, and are more likely to have reduced longevity from cardiovascular causes.
Left-handers may be more likely to suffer bone fractures.
Left-handers have a lower prevalence of arthritis and ulcer.
One systematic review concluded: "Left-handers showed no systematic tendency to suffer from disorders of the immune system".
As handedness is a highly heritable trait associated with various medical conditions, and because many of these conditions could have presented a Darwinian fitness challenge in ancestral populations, this indicates left-handedness may have previously been rarer than it currently is, due to natural selection. However, on average, left-handers have been found to have an advantage in fighting and competitive, interactive sports, which could have increased their reproductive success in ancestral populations.
Income
In 2006, researchers from Lafayette College and Johns Hopkins University concluded that there was no statistically significant correlation between handedness and earnings for the general population, but among college-educated people, left-handers earned 10 to 15% more than their right-handed counterparts.
In a 2014 study published by the National Bureau of Economic Research, Harvard economist Joshua Goodman finds that left-handed people earn 10 to 12 percent less over the course of their lives than right-handed people. Goodman attributes this disparity to higher rates of emotional and behavioral problems in left-handed people.
Sports
Interactive sports such as table tennis, badminton and cricket have an overrepresentation of left-handedness, while non-interactive sports such as swimming show no overrepresentation. Smaller physical distance between participants increases the overrepresentation. In fencing, about half the participants are left-handed. In tennis, 40% of the seeded players are left-handed. The term southpaw is sometimes used to refer to a left-handed individual, especially in baseball and boxing. Some studies suggest that right handed male athletes tend to be statistically taller and heavier than left handed ones.
Other, sports-specific factors may increase or decrease the advantage left-handers usually hold in one-on-one situations:
In cricket, the overall advantage of a bowler's left-handedness exceeds that resulting from experience alone: even disregarding the experience factor (i.e., even for a batter whose experience against left-handed bowlers equals their experience against right-handed bowlers), a left-handed bowler challenges the average (i.e., right-handed) batter more than a right-handed bowler does, because the angle of a bowler's delivery to an opposite-handed batter is much more penetrating than that of a bowler to a same-handed batter (see Wasim Akram).
In baseball, a right-handed pitcher's curve ball will break away from a right-handed batter and towards a left-handed batter (batting left or right does not indicate left or right handedness). While studies of handedness show that only 10% of the general population is left-handed, the proportion of left-handed MLB players is closer to 39% of hitters and 28% of pitchers, according to 2012 data. Historical batting averages show that left-handed batters have a slight advantage over right-handed batters when facing right-handed pitchers. Because there are fewer left-handed pitchers than right-handed pitchers, left-handed batters have more opportunities to face right-handed pitchers than their right-handed counterparts have against left-handed pitchers. Fifteen of the top twenty career batting average leaders in Major League Baseball history have been posted by left-handed batters. Left-handed batters have a slightly shorter run from the batter's box to first base than right-handers. This gives left-handers a slight advantage in beating throws to first base on infield ground balls. Perhaps more important, the follow through of a left-handed swing provides momentum in the direction of first base, while the right handed batter must overcome the swing momentum towards third base before beginning his run.
Because a left-handed pitcher faces first base when he is in position to throw to the batter, whereas a right-handed pitcher has his back to first base, a left-handed pitcher has an advantage when attempting to pick off baserunners at first base.
Defensively in baseball, left-handedness is considered an advantage for first basemen because they are better suited to fielding balls hit in the gap between first and second base, and because they do not have to pivot their body around before throwing the ball to another infielder. For the same reason, the other infielder's positions are seen as being advantageous to right-handed throwers. Historically, there have been few left-handed catchers because of the perceived disadvantage a left-handed catcher would have in making the throw to third base, especially with a right-handed hitter at the plate. A left-handed catcher would have a potentially more dangerous time tagging out a baserunner trying to score. With the ball in the glove on the right hand, a left-handed catcher would have to turn his body to the left to tag a runner. In doing so, he can lose the opportunity to brace himself for an impending collision. On the other hand, the Encyclopedia of Baseball Catchers states:
In four wall handball, typical strategy is to play along the left wall forcing the opponent to use their left hand to counter the attack and playing into the strength of a left-handed competitor.
In handball, left-handed players have an advantage on the right side of the field when attacking, getting a better angle, and that defenders might be unused to them. Since few people are left-handed, there is a demand for such players.
In water polo, the centre forward position has an advantage in turning to shoot on net when rotating the reverse direction as expected by the centre of the opposition defence and gain an improved position to score. Left-handed drivers are usually on the right side of the field, because they can get better angles to pass the ball or shoot for goal.
Ice hockey typically uses a strategy in which a defence pairing includes one left-handed and one right-handed defender. A disproportionately large number of ice hockey players of all positions, 62 percent, shoot left, though this does not necessarily indicate left-handedness.
In American football, the handedness of a quarterback affects blocking patterns on the offensive line. Tight ends, when only one is used, typically line up on the same side as the throwing hand of the quarterback, while the offensive tackle on the opposite hand, which protects the quarterback's "blind side", is typically the most valued member of the offensive line. Receivers also have to adapt to the opposite spin. While uncommon, there have been several notable left-handed quarterbacks.
In bowling, the oil pattern used on the bowling lane breaks down faster the more times a ball is rolled down the lane. Bowlers must continually adjust their shots to compensate for the ball's change in rotation as the game or series is played and the oil is altered from its original pattern. A left-handed bowler competes on the opposite side of the lane from the right-handed bowler and therefore deals with less breakdown of the original oil placement. This means left-handed bowlers have to adjust their shot less frequently than right-handed bowlers in team events or qualifying rounds where there are possibly 4-10 people per set of two lanes. This can allow them to stay more consistent. However, this advantage is not present in bracket rounds and tournament finals where matches are 1v1 on a pair of lanes.
Sex
According to a meta-analysis of 144 studies, totaling 1,787,629 participants, the best estimate for the male to female odds ratio was 1.23, indicating that men are 23% more likely to be left-handed. For example, if the incidence of female left-handedness was 10%, then the incidence of male left-handedness would be approximately 12% (10% incidence of left-handedness among women multiplied by an odds ratio of 1:1.23 for women:men results in a 12.3% incidence of left-handedness among men).
Sexuality and gender identity
Some studies examining the relationship between handedness and sexual orientation have reported that a disproportionate minority of homosexual people exhibit left-handedness, though findings are mixed.
A 2001 study also found that people assigned male at birth whose gender identity did not align with their assigned sex, were more than twice as likely to be left-handed than a clinical control group (19.5% vs. 8.3%, respectively).
Paraphilias (atypical sexual interests) have also been linked to higher rates of left-handedness. A 2008 study analyzing the sexual fantasies of 200 males found "elevated paraphilic interests were correlated with elevated non-right handedness". Greater rates of left-handedness have also been documented among pedophiles.
A 2014 study attempting to analyze the biological markers of asexuality asserts that non-sexual men and women were 2.4 and 2.5 times, respectively, more likely to be left-handed than their heterosexual counterparts.
Mortality rates in combat
A study at Durham University—which examined mortality data for cricketers whose handedness was a matter of public record—found that left-handed men were almost twice as likely to die in war as their right-handed contemporaries. The study theorised that this was because weapons and other equipment was designed for the right-handed. "I can sympathise with all those left-handed cricketers who have gone to an early grave trying desperately to shoot straight with a right-handed Lee Enfield .303", wrote a journalist reviewing the study in the cricket press. The findings echo those of previous American studies, which found that left-handed US sailors were 34% more likely to have a serious accident than their right-handed counterparts.
Episodic memory
A high level of handedness (whether strongly favoring right or left) is associated with poorer episodic memory, and with poorer communication between brain hemispheres, which may give poorer emotional processing, although bilateral stimulation may reduce such effects.
Corpus callosum
A high level of handedness is associated with a smaller corpus callosum whereas low handedness with a larger one.
Divergent thinking
Left-handedness is associated with better divergent thinking.
Products for left-handed use
Many tools and procedures are designed to facilitate use by right-handed people, often without realizing the difficulties incurred by the left-handed. John W. Santrock has written, "For centuries, left-handers have suffered unfair discrimination in a world designed for right-handers."
Many products for left-handed use are made by specialist producers, although not available from normal suppliers. Items as simple as a knife ground for use with the right hand are less convenient for left-handers. There is a multitude of examples: kitchen tools such as knives, corkscrews and scissors, garden tools, and so on. While not requiring a purpose-designed product, there are more appropriate ways for left-handers to tie shoelaces. There are companies that supply products designed specifically for left-handed use. One such is Anything Left-Handed, which in 1967 opened a shop in Soho, London; the shop closed in 2006, but the company continues to supply left-handed products worldwide by mail order.
Writing from left to right as in many languages, in particular, with the left hand covers and tends to smear (depending upon ink drying) what was just written. Left-handed writers have developed various ways of holding a pen for best results. For using a fountain pen, preferred by many left-handers, nibs ground to optimise left-handed use (pushing rather than pulling across the paper) without scratching are available.
Bias against left-handers
McManus noted that, as the Industrial Revolution spread across Western Europe and the United States in the 19th century, workers needed to operate complex machines that were designed with right-handers in mind. This would have made left-handers more visible and at the same time appear less capable and more clumsy. Writing left-handed with a dip pen, in particular, was prone to blots and smearing.
Negative connotations and discrimination
Moreover, apart from inconvenience, left-handed people have historically been considered unlucky or even malicious for their difference by the right-handed majority. In many languages, including English, the word for the direction "right" also means "correct" or "proper". Throughout history, being left-handed was considered negative, or evil.
The Latin adjective means as well as , and this double meaning survives in European derivatives of Latin, including the English words sinister (meaning both 'evil' and 'on the bearer's left on a coat of arms') and ambisinister meaning 'awkward or clumsy with both or either hand'.
There are many negative connotations associated with the phrase left-handed: clumsy, awkward, unlucky, insincere, sinister, malicious, and so on. A "left-handed compliment" is one that has two meanings, one of which is unflattering to the recipient. In French, means both and or , while (cognate to English direct and related to adroit) means both and , as well as and the legal sense of . The name Dexter derives from the Latin for , as does the word dexterity meaning manual skill. As these are all very old words, they would tend to support theories indicating that the predominance of right-handedness is an extremely old phenomenon.
Black magic is sometimes referred to as the "left-hand path".
Discrimination in education
Before the development of fountain pens and other writing instruments, children were taught to write with a dip pen. While a right-hander could smoothly drag the pen across paper from left to right, a dip pen could not easily be pushed across by the left hand without digging into the paper and making blots and stains. Even with more modern pens, writing from left to right, as in many languages, with the left hand covers and can smear what was just written when moving across the line.
Into the 20th and even the 21st century, left-handed children in Uganda were beaten by schoolteachers or parents for writing with their left hand, or had their left hands tied behind their backs to force them to write with their right hand. As a child, the future British king George VI (1895–1952) was naturally left-handed. He was forced to write with his right hand, as was common practice at the time. He was not expected to become king, so that was not a factor.
Until very recently in Taiwan, left-handed people were forced to switch to being right-handed, or at least switch to writing with the right hand. Due to the importance of stroke order, developed for the comfortable use of right-handed people, it is considered more difficult to write legible Chinese characters with the left hand than it is to write Latin letters, though difficulty is subjective and depends on the writer. Because writing when moving one's hand away from its side towards the other side of the body can cause smudging if the outward side of the hand is allowed to drag across the writing, writing in the Latin alphabet might possibly be less feasible with the left hand than the right under certain circumstances. Conversely, right-to-left alphabets, such as the Arabic and Hebrew, are generally considered easier to write with the left hand. Depending on the position and inclination of the writing paper, and the writing method, the left-handed writer can write as neatly and efficiently or as messily and slowly as right-handed writers. Usually the left-handed child needs to be taught how to write correctly with the left hand, since discovering a comfortable left-handed writing method on one's own may not be straightforward.
In the Soviet school system, all left-handed children were forced to write with their right hand.
International Left-Handers Day
International Left-Handers Day is held annually every August 13. It was founded by the Left-Handers Club in 1992, with the club itself having been founded in 1990. International Left-Handers Day is, according to the club, "an annual event when left-handers everywhere can celebrate their sinistrality (left-handedness) and increase public awareness of the advantages and disadvantages of being left-handed." It celebrates their uniqueness and differences, who are from seven to ten percent of the world's population. Thousands of left-handed people in today's society have to adapt to use right-handed tools and objects. Again according to the club, "in the U.K. alone there were over 20 regional events to mark the day in 2001 – including left-v-right sports matches, a left-handed tea party, pubs using left-handed corkscrews where patrons drank and played pub games with the left hand only, and nationwide 'Lefty Zones' where left-handers' creativity, adaptability and sporting prowess were celebrated, whilst right-handers were encouraged to try out everyday left-handed objects to see just how awkward it can feel using the wrong equipment."
In other animals
Kangaroos and other macropod marsupials show a left-hand preference for everyday tasks in the wild. 'True' handedness is unexpected in marsupials however, because unlike placental mammals, they lack a corpus callosum. Left-handedness was particularly apparent in the red kangaroo (Macropus rufus) and the eastern gray kangaroo (Macropus giganteus). Red-necked (Bennett's) wallabies (Macropus rufogriseus) preferentially use their left hand for behaviours that involve fine manipulation, but the right for behaviours that require more physical strength. There was less evidence for handedness in arboreal species. Studies of dogs, horses, and domestic cats have shown that females of those species tend to be right-handed, while males tend to be left-handed.
| Biology and health sciences | Human anatomy | Health |
172716 | https://en.wikipedia.org/wiki/Caltrop | Caltrop | A caltrop (also known as caltrap, galtrop, cheval trap, galthrap, galtrap, calthrop, jackrock or crow's foot) is an area denial weapon made up of usually four, but possibly more, sharp nails or spines arranged in such a manner that one of them always points upward from a stable base (for example, a tetrahedron). Historically, caltrops were part of defences that served to slow the advance of troops, especially horses, chariots, and war elephants, and were particularly effective against the soft feet of camels. In modern times, caltrops are effective when used against wheeled vehicles with pneumatic tires.
Name
The modern name "caltrop" is derived from the Old English (heel-trap), such as in the French usage (shoe-trap). The Latin word originally referred to this and provides part of the modern scientific name of a plant commonly called the caltrop, Tribulus terrestris, whose spiked seed cases resemble caltrops and can injure feet and puncture bicycle tires. This plant can also be compared to Centaurea calcitrapa, which is also sometimes referred to as the "caltrop". Trapa natans, a water plant with similarly shaped spiked seeds and edible fruit, is called the "water caltrop".
History
The caltrop was called by the ancient Romans, or sometimes , the latter meaning "jagged iron" (literally "iron spiny snail-shell"). The former term derives from the ancient Greek word meaning three spikes.
The late Roman writer Vegetius, referring in his work De re militari to scythed chariots, wrote:
Another example of the use of caltrops was found in Jamestown, Virginia, in the United States:
The Japanese version of the caltrop is called . Makibishi were sharp spiked objects that were used in feudal Japan to slow pursuers and also were used in the defence of samurai fortifications. Iron makibishi were called , while the makibishi made from the dried seed pod of the water caltrop, or water chestnut (genus Trapa), formed a natural type of makibashi called . Both types of makibishi could penetrate the thin soles of shoes, such as the sandals, which were commonly worn in feudal Japan.
Modern uses
World War I
During service in World War I, Australian Light Horse troops collected caltrops as keepsakes. These caltrops were either made by welding two pieces of wire together to form a four-pointed star or pouring molten steel into a mould to form a solid, seven-pointed star. The purpose of these devices was to disable horses. They were exchanged with French troops for bullets. The Australian Light Horse troops referred to them as "Horse Chestnuts".
World War II
Caltrops were used extensively and effectively during World War II. The modifications and variants produced by the Special Operations Executive (SOE) and the Office of Strategic Services (OSS) of the United States are still in use today within special forces and law enforcement bodies.
The Germans dropped crow's feet (). These were made from two segments of sheet metal welded together into a tetrapod with four barbed points and then painted in camouflage colours. They came in two sizes with a side length of either . They were dropped from aircraft in containers the same size as bombs and were dispersed by a small explosive charge.
Tire deflation device
Inventors patented caltrop-like devices to deflate vehicle tires in a manner useful to law enforcement agencies or the military. They are currently used by the military and police.
Labour disputes
Caltrops have been used at times during labour strikes and other disputes. Such devices were used by some to destroy the tires of management and replacement workers.
Caltrops, referred to as "jack rocks" in news articles, were used during the Caterpillar strike in 1995, puncturing tires on vehicles crossing the picket line in Peoria, Illinois. Because of their small size and the difficulty proving their source, both the company and the United Auto Workers blamed each other. Collateral damage included a school bus and a walking mail carrier. In Illinois, the state legislature passed a law making the possession of such devices a misdemeanor.
Via drones
During the Russian invasion of Ukraine, Ukraine has used drones to drop caltrops on key roads to disrupt wheeled vehicles carrying Russian military materiel, and make them easier to target with loitering munitions.
Symbol
A caltrop has a variety of symbolic uses and is commonly found as a charge in heraldry. For instance, the Finnish noble family Fotangel (Swedish for 'caltrop') had arms gules, three caltrops argent.
It has also been adopted by military units: the caltrop is the symbol of the US Army's III Corps, which is based at Fort Cavazos, Texas. III Corps traces its lineage to the days of horse cavalry, which used the caltrop as a defensive area denial weapon.
The caltrop is also the symbol of the United States Marine Corps' 3rd Division, formed on 16 September 1942.
Similar devices
Punji sticks perform a similar role to caltrops. These are sharpened sticks placed vertically in the ground. Their use in modern times targets the body and limbs of a falling victim by means of a pit or tripwire.
During the Second World War, large caltrop-shaped objects made from reinforced concrete were used as anti-tank devices, although it seems that these were rare. Much more common were concrete devices called dragon's teeth, which were designed to wedge into tank treads. Large ones weighing over are still used defensively to deny access to wheeled vehicles, especially in camp areas. As dragon's teeth are immobile, the analogy with the caltrop is inexact. Another caltrop-like defence during World War II was the massive steel, freestanding Czech hedgehog; the works were designed as anti-tank obstacles and could also damage landing craft and warships that came too close to shore. These were used by the Germans to defend beaches in Normandy and other coastal areas. Czech hedgehogs are heavily featured and plainly visible in the 1998 Steven Spielberg-directed American epic war film Saving Private Ryan, throughout the scenes early in the film depicting the June6, 1944 Omaha Beach assault (part of the Normandy landings during World War II).
Tetrapods are concrete blocks shaped like caltrops, which interlock when piled up. They are used as riprap in the construction of breakwaters and other sea defences, as they have been found to let the water pass through them and interrupt natural processes less than some other defenses.
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