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3183196	https://en.wikipedia.org/wiki/Rolling%20%28metalworking%29	Rolling (metalworking)	In metalworking, rolling is a metal forming process in which metal stock is passed through one or more pairs of rolls to reduce the thickness, to make the thickness uniform, and/or to impart a desired mechanical property. The concept is similar to the rolling of dough. Rolling is classified according to the temperature of the metal rolled. If the temperature of the metal is above its recrystallization temperature, then the process is known as hot rolling. If the temperature of the metal is below its recrystallization temperature, the process is known as cold rolling. In terms of usage, hot rolling processes more tonnage than any other manufacturing process, and cold rolling processes the most tonnage out of all cold working processes. Roll stands holding pairs of rolls are grouped together into rolling mills that can quickly process metal, typically steel, into products such as structural steel (I-beams, angle stock, channel stock), bar stock, and rails. Most steel mills have rolling mill divisions that convert the semi-finished casting products into finished products. There are many types of rolling processes, including ring rolling, roll bending, roll forming, profile rolling, and controlled rolling. Iron and steel The earliest rolling mills in crude form but the same basic principles were found in Middle East and South Asia as early as 600 BCE. The invention of the rolling mill in Europe may be attributed to Leonardo da Vinci in his drawings. Earliest rolling mills were slitting mills, which were introduced from what is now Belgium to England in 1590. These passed flat bars between rolls to form a plate of iron, which was then passed between grooved rolls (slitters) to produce rods of iron. The first experiments at rolling iron for tinplate took place about 1670. In 1697, Major John Hanbury erected a mill at Pontypool to roll "Pontypool plates" – blackplate. Later this began to be rerolled and tinned to make tinplate. The earlier production of plate iron in Europe had been in forges, not rolling mills. The slitting mill was adapted to producing hoops (for barrels) and iron with a half-round or other sections by means that were the subject of two patents of c. 1679. Some of the earliest literature on rolling mills can be traced back to the Swedish engineer Christopher Polhem in his Patriotista Testamente of 1761, where he mentions rolling mills for both plate and bar iron. He also explains how rolling mills can save on time and labor because a rolling mill can produce 10 to 20 or more bars at the same time. A patent was granted to Thomas Blockley of England in 1759 for the polishing and rolling of metals. Another patent was granted in 1766 to Richard Ford of England for the first tandem mill. A tandem mill is one in which the metal is rolled in successive stands; Ford's tandem mill was for hot rolling of wire rods. Other metals Rolling mills for lead seem to have existed by the late 17th century. Copper and brass were also rolled by the late 18th century. Modern rolling Until well into the eighteenth century, rolling mills derived their power from water wheels. The first recorded use of a steam engine directly driving a mill is attributed to John Wilkinson's Bradley Works where, in 1786, a Boulton and Watt engine was coupled to a slitting and rolling mill. The use of steam engines considerably enhanced the production capabilities of the mills, until this form of power was displaced by electric motors soon after 1900. Modern rolling practice can be attributed to the pioneering efforts of Henry Cort of Funtley Iron Mills, near Fareham in Hampshire, England. In 1783, a patent number was issued to Henry Cort for his use of grooved rolls for rolling iron bars. With this new design, mills were able to produce 15 times more output per day than with a hammer. Although Cort was not the first to use grooved rolls, he was the first to combine the use of many of the best features of various ironmaking and shaping processes known at the time. Thus modern writers have called him "father of modern rolling". The first rail rolling mill was established by John Birkenshaw at Bedlington Ironworks in Northumberland, England, in 1820, where he produced fish-bellied wrought iron rails in lengths of 15 to 18 feet. With the advancement of technology in rolling mills, the size of rolling mills grew rapidly along with the size of the products being rolled. One example of this was at The Great Exhibition in London in 1851, where a plate 20 feet long, 3 feet wide, and 7/16 of an inch thick, and weighing 1,125 pounds, was exhibited by the Consett Iron Company. Further evolution of the rolling mill came with the introduction of three-high mills in 1853 used for rolling heavy sections. Hot and cold rolling Hot rolling Hot rolling is a metalworking process that occurs above the recrystallization temperature of the material. After the grains deform during processing, they recrystallize, which maintains an equiaxed microstructure and prevents the metal from work hardening. The starting material is usually large pieces of metal, like semi-finished casting products, such as ingots, slabs, blooms, and billets. If these products came from a continuous casting operation, the products are usually fed directly into the rolling mills at the proper temperature. In smaller operations, the material starts at room temperature and must be heated. This is done in a gas- or oil-fired soaking pit for larger workpieces; for smaller workpieces, induction heating is used. As the material is worked, the temperature must be monitored to make sure it remains above the recrystallization temperature. To maintain a safety factor a finishing temperature is defined above the recrystallization temperature; this is usually above the recrystallization temperature. If the temperature does drop below this temperature the material must be re-heated prior to additional hot rolling. Hot-rolled metals generally have little directionality in their mechanical properties or deformation-induced residual stresses. However, in certain instances non-metallic inclusions will impart some directionality and workpieces less than thick often have some directional properties. Non-uniform cooling will induce a lot of residual stresses, which usually occurs in shapes that have a non-uniform cross-section, such as I-beams. While the finished product is of good quality, the surface is covered in mill scale, which is an oxide that forms at high temperatures. It is usually removed via pickling or the smooth clean surface (SCS) process, which reveals a smooth surface. Dimensional tolerances are usually 2 to 5% of the overall dimension. Hot-rolled mild steel seems to have a wider tolerance for the level of included carbon than does cold-rolled steel, and is, therefore, more difficult for a blacksmith to use. Hot rolling is used mainly to produce sheet metal or simple cross-sections, such as rail tracks. Shape rolling design Rolling mills are often divided into roughing, intermediate and finishing rolling cages. During shape rolling, an initial billet (round or square) with edge of diameter typically ranging between 100 and 140 mm is continuously deformed to produce a certain finished product with smaller cross section dimension and geometry. Starting from a given billet, different sequences can be adopted to produce a certain final product. However, since each rolling mill is significantly expensive (up to 2 million euros), a typical requirement is to reduce the number of rolling passes. Different approaches have been achieved, including empirical knowledge, employment of numerical models, and Artificial Intelligence techniques. Lambiase et al. validated a finite element model (FE) for predicting the final shape of a rolled bar in round-flat pass. One of the major concerns when designing rolling mills is to reduce the number of passes. A possible solution to such requirements is the slit pass, also called split pass, which divides an incoming bar in two or more subparts, thus virtually increasing the cross section reduction ratio per pass as reported by Lambiase. Another solution for reducing the number of passes in rolling mills is the employment of automated systems for Roll Pass Design as that proposed by Lambiase and Langella. subsequently, Lambiase further developed an Automated System based on Artificial Intelligence and particularly an integrated system including an inferential engine based on Genetic Algorithms a knowledge database based on an Artificial Neural Network trained by a parametric Finite element model and to optimize and automatically design rolling mills. Cold rolling Cold rolling occurs with the metal below its recrystallization temperature (usually at room temperature), which increases the strength via strain hardening up to 20%. It also improves the surface finish and holds tighter tolerances. Commonly cold-rolled products include sheets, strips, bars, and rods; these products are usually smaller than the same products that are hot rolled. Because of the smaller size of the workpieces and their greater strength, as compared to hot rolled stock, four-high or cluster mills are used. Cold rolling cannot reduce the thickness of a workpiece as much as hot rolling in a single pass. Cold-rolled sheets and strips come in various conditions: full-hard, half-hard, quarter-hard, and skin-rolled. Full-hard rolling reduces the thickness by 50%, while the others involve less of a reduction. Cold rolled steel is then annealed to induce ductility in the cold rolled steel which is simply known as a Cold Rolled and Close Annealed. Skin-rolling, also known as a skin-pass, involves the least amount of reduction: 0.5–1%. It is used to produce a smooth surface, a uniform thickness, and reduce the yield point phenomenon (by preventing Lüders bands from forming in later processing). It locks dislocations at the surface and thereby reduces the possibility of formation of Lüders bands. To avoid the formation of Lüders bands it is necessary to create substantial density of unpinned dislocations in ferrite matrix. It is also used to break up the spangles in galvanized steel. Skin-rolled stock is usually used in subsequent cold-working processes where good ductility is required. Other shapes can be cold-rolled if the cross-section is relatively uniform and the transverse dimension is relatively small. Cold rolling shapes requires a series of shaping operations, usually along the lines of sizing, breakdown, roughing, semi-roughing, semi-finishing, and finishing. If processed by a blacksmith, the smoother, more consistent, and lower levels of carbon encapsulated in the steel makes it easier to process, but at the cost of being more expensive. Processes Roll bending Roll bending produces a cylindrical shaped product from plate or steel metals. Roll forming Roll forming, roll bending or plate rolling is a continuous bending operation in which a long strip of metal (typically coiled steel) is passed through consecutive sets of rolls, or stands, each performing only an incremental part of the bend, until the desired cross-section profile is obtained. Roll forming is ideal for producing parts with long lengths or in large quantities. There are three main processes: 4 rollers, 3 rollers and 2 rollers, each of which has as different advantages according to the desired specifications of the output plate. Flat rolling Flat rolling is the most basic form of rolling with the starting and ending material having a rectangular cross-section. The material is fed in between two rollers, called working rolls, that rotate in opposite directions. The gap between the two rolls is less than the thickness of the starting material, which causes it to deform. The decrease in material thickness causes the material to elongate. The friction at the interface between the material and the rolls causes the material to be pushed through. The amount of deformation possible in a single pass is limited by the friction between the rolls; if the change in thickness is too great the rolls just slip over the material and do not draw it in. The final product is either sheet or plate, with the former being less than thick and the latter greater than; however, heavy plates tend to be formed using a press, which is termed forging, rather than rolling. Often the rolls are heated to assist in the workability of the metal. Lubrication is often used to keep the workpiece from sticking to the rolls. To fine-tune the process, the speed of the rolls and the temperature of the rollers are adjusted. For thin sheet metal with a thickness less than , the rolling is done in a cluster mill because the small thickness requires a small diameter rolls. To reduce the need for small rolls pack rolling is used, which rolls multiple sheets together to increase the effective starting thickness. As the foil sheets come through the rollers, they are trimmed and slitted with circular or razor-like knives. Trimming refers to the edges of the foil, while slitting involves cutting it into several sheets. Aluminum foil is the most commonly produced product via pack rolling. This is evident from the two different surface finishes; the shiny side is on the roll side and the dull side is against the other sheet of foil. Ring rolling Ring rolling is a specialized type of hot rolling that increases the diameter of a ring. The starting material is a thick-walled ring. This workpiece is placed between two rolls, an inner idler roll and a driven roll, which presses the ring from the outside. As the rolling occurs the wall thickness decreases as the diameter increases. The rolls may be shaped to form various cross-sectional shapes. The resulting grain structure is circumferential, which gives better mechanical properties. Diameters can be as large as and face heights as tall as . Common applications include railway tyres, bearings, gears, rockets, turbines, airplanes, pipes, and pressure vessels. Structural shape rolling Controlled rolling Controlled rolling is a type of thermomechanical processing which integrates controlled deformation and heat treating. The heat which brings the workpiece above the recrystallization temperature is also used to perform the heat treatments so that any subsequent heat treating is unnecessary. Types of heat treatments include the production of a fine grain structure; controlling the nature, size, and distribution of various transformation products (such as ferrite, austenite, pearlite, bainite, and martensite in steel); inducing precipitation hardening; and, controlling the toughness. In order to achieve this the entire process must be closely monitored and controlled. Common variables in controlled rolling include the starting material composition and structure, deformation levels, temperatures at various stages, and cool-down conditions. The benefits of controlled rolling include better mechanical properties and energy savings. Forge rolling Forge rolling is a longitudinal rolling process to reduce the cross-sectional area of heated bars or billets by leading them between two contrary rotating roll segments. The process is mainly used to provide optimized material distribution for subsequent die forging processes. Owing to this a better material utilization, lower process forces and better surface quality of parts can be achieved in die forging processes. Basically any forgeable metal can also be forge-rolled. Forge rolling is mainly used to preform long-scaled billets through targeted mass distribution for parts such as crankshafts, connection rods, steering knuckles and vehicle axles. Narrowest manufacturing tolerances can only partially be achieved by forge rolling. This is the main reason why forge rolling is rarely used for finishing, but mainly for preforming. Characteristics of forge rolling: high productivity and high material utilization good surface quality of forge-rolled workpieces extended tool life-time small tools and low tool costs improved mechanical properties due to optimized grain flow compared to exclusively die forged workpieces Mills A rolling mill, also known as a reduction mill or mill, has a common construction independent of the specific type of rolling being performed: Work rolls Backup rolls – are intended to provide rigid support required by the working rolls to prevent bending under the rolling load Rolling balance system – to ensure that the upper work and back up rolls are maintained in proper position relative to lower rolls Roll changing devices – use of an overhead crane and a unit designed to attach to the neck of the roll to be removed from or inserted into the mill. Mill protection devices – to ensure that forces applied to the backup roll chocks are not of such a magnitude to fracture the roll necks or damage the mill housing Roll cooling and lubrication systems Pinions – gears to divide power between the two spindles, rotating them at the same speed but in different directions Gearing – to establish desired rolling speed Drive motors – rolling narrow foil product to thousands of horsepower Electrical controls – constant and variable voltages applied to the motors Coilers and uncoilers – to unroll and roll up coils of metal Slabs are the feed material for hot strip mills or plate mills and blooms are rolled to billets in a billet mill or large sections in a structural mill. The output from a strip mill is coiled and, subsequently, used as the feed for a cold rolling mill or used directly by fabricators. Billets, for re-rolling, are subsequently rolled in either a merchant, bar or rod mill. Merchant or bar mills produce a variety of shaped products such as angles, channels, beams, rounds (long or coiled) and hexagons. Configurations Mills are designed in different types of configurations, with the most basic being a two-high non-reversing, which means there are two rolls that only turn in one direction. The two-high reversing mill has rolls that can rotate in both directions, but the disadvantage is that the rolls must be stopped, reversed, and then brought back up to rolling speed between each pass. To resolve this, the three-high mill was invented, which uses three rolls that rotate in one direction; the metal is fed through two of the rolls and then returned through the other pair. The disadvantage to this system is the workpiece must be lifted and lowered using an elevator. All of these mills are usually used for primary rolling and the roll diameters range from . To minimize the roll diameter a four-high or cluster mill is used. A small roll diameter is advantageous because less roll is in contact with the material, which results in a lower force and power requirement. The problem with a small roll is a reduction of stiffness, which is overcome using backup rolls. These backup rolls are larger and contact the back side of the smaller rolls. A four-high mill has four rolls, two small and two large. A cluster mill has more than four rolls, usually in three tiers. These types of mills are commonly used to hot roll wide plates, most cold rolling applications, and to roll foils. Historically mills were classified by the product produced: Blooming, cogging and slabbing mills, being the preparatory mills to rolling finished rails, shapes or plates, respectively. If reversing, they are from 34 to 48 inches in diameter, and if three-high, from 28 to 42 inches in diameter. Billet mills, three-high, rolls from 24 to 32 inches in diameter, used for the further reduction of blooms down to 1.5x1.5-inch billets, being the nubpreparatory mills for the bar and rod Beam mills, three-high, rolls from 28 to 36 inches in diameter, for the production of heavy beams and channels 12 inches and over. Rail mills with rolls from 26 to 40 inches in diameter. Shape mills with rolls from 20 to 26 inches in diameter, for smaller sizes of beams and channels and other structural shapes. Merchant bar mills with rolls from 16 to 20 inches in diameter. Small merchant bar mills with finishing rolls from 8 to 16 inches in diameter, generally arranged with a larger size roughing stand. Rod and wire mills with finishing rolls from 8 to 12 inches in diameter, always arranged with larger size roughing stands. Hoop and cotton tie mills, similar to small merchant bar mills. Armour plate mills with rolls from 44 to 50 inches in diameter and 140 to 180-inch body. Plate mills with rolls from 28 to 44 inches in diameter. Sheet mills with rolls from 20 to 32 inches in diameter. Universal mills for the production of square-edged or so-called universal plates and various wide flanged shapes by a system of vertical and horizontal rolls. Tandem mill A tandem mill is a special type of modern rolling mill where rolling is done in one pass. In a traditional rolling mill rolling is done in several passes, but in tandem mill there are several stands (>=2 stands) and reductions take place successively. The number of stands ranges from 2 to 18. Tandem mills can be either of hot or cold rolling mill types. Cold rolling mills may be further divided into continuous or batch processing. A continuous mill has a looping tower which allows the mill to continue rolling slowly the strip in the tower, while a strip welder joins the tail of the current coil to the head of the next coil. At the exit end of the mill there is normally a flying shear (to cut the strip at or near the weld) followed by two coilers; one being unloaded while the other winds on the current coil. Looping towers are also used in other places; such as continuous annealing lines and continuous electrolytic tinning and continuous galvanising lines. Defects Thickness changes along length In hot rolling, if the temperature of the workpiece is not uniform the flow of the material will occur more in the warmer parts and less in the cooler. If the temperature difference is great enough cracking and tearing can occur. The cooler sections are, among other things, a result of the supports in the re-heat furnace. When cold rolling, virtually all of the strip thickness variation is the result of the eccentricity and out-of-roundness of the Back-up Rolls from about Stand 3 of the Hot Strip Mill through to the Finished Product. The Back-up Roll eccentricity can be up to 100 μm in magnitude per stack. The eccentricity can be measured off-line by plotting the force variation against time with the Mill on creep, no strip present, and the Mill Stand below face. A modified Fourier analysis was employed by the 5 Stand Cold Mill at Bluescope Steel, Port Kembla from 1986 until that Cold Mill ceased production in 2009. Within each coil, the exit thickness deviation times 10 for every meter of strip was stored in a file. This file was analyzed separately for each frequency/wavelength from 5 m to 60 m in steps of 0.1 m. To improve the accuracy, care was taken to use a full multiple of each wavelength (100*). The calculate amplitudes were plotted against the wavelength, so that the spikes could be compared to the expected wavelengths created by the Backup Rolls of each Stand. If a Mill Stand is fitted with Hydraulic Pistons in series with, or instead of the electrically driven Mechanical Screws, then it is possible to eliminate the effect of that Stands Back-up Roll eccentricity. While rolling, the eccentricity of each Back-up Roll is determined by sampling the roll force and assigning it to the corresponding portion of each Back-up Roll's rotational position. These recordings are then used to operate the Hydraulic Piston so as to neutralize the eccentricities. Flatness and shape In a flat metal workpiece, the flatness is a descriptive attribute characterizing the extent of the geometric deviation from a reference plane. The deviation from complete flatness is the direct result of the workpiece relaxation after hot or cold rolling, due to the internal stress pattern caused by the non-uniform transversal compressive action of the rolls and the uneven geometrical properties of the entry material. The transverse distribution of differential strain/elongation-induced stress with respect to the material's average applied stress is commonly referenced to as shape. Due to the strict relationship between shape and flatness, these terms can be used in an interchangeable manner. In the case of metal strips and sheets, the flatness reflects the differential fiber elongation across the width of the workpiece. This property must be subject to an accurate feedback-based control in order to guarantee the machinability of the metal sheets in the final transformation processes. Some technological details about the feedback control of flatness are given in. Profile Profile is made up of the measurements of crown and wedge. Crown is the thickness in the center as compared to the average thickness at the edges of the workpiece. Wedge is a measure of the thickness at one edge as opposed to the other edge. Both may be expressed as absolute measurements or as relative measurements. For instance, one could have 2 mil of crown (the center of the workpiece is 2 mil thicker than the edges), or one could have 2% crown (the center of the workpiece is 2% thicker than the edges). It is typically desirable to have some crown in the workpiece as this will cause the workpiece to tend to pull to the center of the mill, and thus will run with higher stability. Flatness Maintaining a uniform gap between the rolls is difficult because the rolls deflect under the load required to deform the workpiece. The deflection causes the workpiece to be thinner on the edges and thicker in the middle. This can be overcome by using a crowned roller (parabolic crown), however the crowned roller will only compensate for one set of conditions, specifically the material, temperature, and amount of deformation. Other methods of compensating for roll deformation include continual varying crown (CVC), pair cross rolling, and work roll bending. CVC was developed by SMS-Siemag AG and involves grinding a third order polynomial curve into the work rolls and then shifting the work rolls laterally, equally, and opposite to each other. The effect is that the rolls will have a gap between them that is parabolic in shape, and will vary with lateral shift, thus allowing for control of the crown of the rolls dynamically. Pair cross rolling involves using either flat or parabolically crowned rolls, but shifting the ends at an angle so that the gap between the edges of the rolls will increase or decrease, thus allowing for dynamic crown control. Work roll bending involves using hydraulic cylinders at the ends of the rolls to counteract roll deflection. Another way to overcome deflection issues is by decreasing the load on the rolls, which can be done by applying a longitudinal force; this is essentially drawing. Other method of decreasing roll deflection include increasing the elastic modulus of the roll material and adding back-up supports to the rolls. The different classifications for flatness defects are: Symmetrical edge wave - the edges on both sides of the workpiece are "wavy" due to the material at the edges being longer than the material in the center. Asymmetrical edge wave - one edge is "wavy" due to the material at one side being longer than the other side. Center buckle - The center of the strip is "wavy" due to the strip in the center being longer than the strip at the edges. Quarter buckle - This is a rare defect where the fibers are elongated in the quarter regions (the portion of the strip between the center and the edge). This is normally attributed to using excessive roll bending force since the bending force may not compensate for the roll deflection across the entire length of the roll. One could have a flatness defect even with the workpiece having the same thickness across the width. Also, one could have fairly high crown or wedge, but still produce material that is flat. In order to produce flat material, the material must be reduced by the same percentage across the width. This is important because mass flow of the material must be preserved, and the more a material is reduced, the more it is elongated. If a material is elongated in the same manner across the width, then the flatness coming into the mill will be preserved at the exit of the mill. Draught The difference between the thickness of initial and rolled metal piece is called Draught. Thus if is initial thickness and is final thickness, then the draught is given by The maximum draught that can be achieved via rollers of radius with coefficient of static friction between the roller and the metal surface is given by This is the case when the frictional force on the metal from inlet contact matches the negative force from the exit contact. Surface defect types There are six types of surface defects: Lap This type of defect occurs when a corner or fin is folded over and rolled but not welded into the metal. They appear as seams across the surface of the metal. Mill-shearing These defects occur as a feather-like lap. Rolled-in scale This occurs when mill scale is rolled into metal. Scabs These are long patches of loose metal that have been rolled into the surface of the metal. Seams They are open, broken lines that run along the length of the metal and caused by the presence of scale as well as due to pass roughness of Roughing mill. Slivers Prominent surface ruptures. Surface defect remediation Many surface defects can be scarfed off the surface of semi-finished rolled products before further rolling. Methods of scarfing have included hand-chipping with chisels (18th and 19th centuries); powered chipping and grinding with air chisels and grinders; burning with an oxy-fuel torch, whose gas pressure blows away the metal or slag melted by the flame; and laser scarfing.	Technology	Metallurgy	null
3183570	https://en.wikipedia.org/wiki/Malus%20sieversii	Malus sieversii	Malus sieversii is a wild apple. According to DNA analysis conducted in 2010, it is the primary ancestor of the domesticated apple, M. domestica. Native to Central Asia, M. sieversii prefers warm and damp habitats. Its conservation status is vulnerable. Description It is a deciduous tree growing , very similar in appearance to the domestic apple. Its pollen grains vary in size and are seen to be ovular when dry and spherical when swelled with water. Its fruit is the largest of any species of Malus except domestica, up to 7 cm in diameter, equal in size to many modern apple cultivars. Unlike domesticated varieties, its leaves go red in autumn: 62% of the trees in the wild do this compared to only 2.8% of the regular apple plant or the 2,170 English cultivated varieties. M. sieversii has the capability to reproduce vegetatively as they form root suckers, or basal shoots. The clonal individual grows from the adventitious bud on the root, with identical genetic materials to the mother plant. It was originally believed that wild apples produce root suckers only when upper plant parts are damaged, but further evidence suggests root-sucker growth occurs in healthy plants as a dispersal aid. For wild apples, proper development of root suckers requires certain humidity and aeration levels in the surface soil, where the mother root is located. Successful root-sucker growth also depends on shoot arrangement, time of growth and health conditions of the mother plant. Genetics of self-incompatibility, the system for preventing self-fertilization in angiosperms, have also been studied for M. sieversii. Its genetic diversity in relation to self-incompatibility is substantially less when compared to its close relative, M. sylvestris. Although M. sieversii lacks this diversity, it can survive in the wild without intervention as long as no more diversity loss occurs. The leading theory for this lack of diversity is due to a major population bottleneck during the Last Glacial Maximum which caused wild M. sieversii populations to scale back into a smaller area within the valley of the Ili River. Growth cycle The growth cycle of M. sieversii could be divided into several stages from germination to developing fruit bearing trees, and to the death of aged trees. Period I starts from seed germination to the development of vegetative tree parts, and to the first round of fruiting. A typical apple tree reaches age 6 to 8 in this period. Prominent primary growth and a significant number of root sucker growth are seen in this period. Period II involves more growth and fruit bearing. Wild apple trees usually reach age 10 to 12 in this period. The number of fruits increase significantly as prominent secondary growth and branching take place. Period III involves more growth and fruit bearing. Wild apple trees enter regular fruiting and reach maximum fruit bearing in this period, and usually reach age 25 to 30. Decreased branching rate, and less growth of crown periphery are seen. A lot of vegetative growth is contributed by basal shoots (root suckers). Period IV is fully fruit bearing. Development of trees fully ceases in this stage and they reach around 60 to 70 years of age. Branching decreases as they slowly die off from the base to the periphery. After age 60 to 70, trees experience rapid desiccation of the branches, then death. Taxonomy The species was first described as Pyrus sieversii due to its similarities with pears in 1833 by Carl Friedrich von Ledebour, a German naturalist who saw them growing in the Altai Mountains. Progenitor of cultivated apples Malus sieversii has previously been identified as the main contributor to the genome of the cultivated apple (Malus domestica), on the basis of morphological, molecular, and historical evidence. Fruit traits including crispness, more flavour intensity and fruit weight have undergone differential selection by humans to produce Malus domestica as seen today. The dispersal of M. sieversii and its progeny throughout history can be attributed to the Silk Road. A DNA analysis in 2010 confirmed M. sieversii as the progenitor of the cultivated apple. It has a highly variable genetic diversity therefore it is the genetic source for abiotic and biotic stress tolerance, many disease resistance and unique fruit traits. Planting cultivated apple varieties close to wild groves causes crossbreeding. Distribution and habitat Malus sieversii is distributed mainly within the Ili valley in southern Kazakhstan, where the damp climate suits its growth. It appears in many different habitats. Although the species prefers high temperatures and short winters, it is also found in the Tian Shan mountains which have long and harsh winters. Conservation After the collapse of the USSR and the closure of the Gardening Development Program, the local population began to actively cut down wild gardens in the Zailiyskiy Alatau. The vacated territories are used for building houses and grazing animals. Malus sieversii has been designated as second conservation priority in the China Plant Red Data Book, and has been marked as vulnerable by the International Union for Conservation of Nature (IUCN). Human activities and natural disasters are the major contributors to the decline of M. sieversii natural population. Fungal pathogens, such as Phytophthora plurivora and Alternaria alternata, also play a major role in the decline of M. sieversii populations, by degrading vegetative parts such as the fine root systems. This immune vulnerability makes M. sieversii become susceptible to more parasites, such as pathogenic insects Agrilus mali, to further destroy the population. Ex situ conservation, or seed banking, is believed to be a feasible long-term resolution to protect its genetic diversity, and has been seen in the United States using seeds collected from Kazakhstan and the Kyrgyz Republic. In situ conservation was also found with barbed wire fences being placed around regions distributed with M. sieversii, as seen in areas within Xinyuan, China. A study in 2016 has shown the effectiveness of protecting M. sieversii populations in situ through stratification and seed coat removal. Also, in situ enclosures are more effective in higher elevations as they are at less risk of human and insect injuries. Aside from traditional conservation methods, biofertilizer has shown effective results inhibiting fungal pathogen, Alternaria alternata in wild apple trees. It does this by improving antioxidant capability of wild apple trees following the infection, promoting root growth and enhancing soil metabolism. Recently, a combination of innovative methods including cloning and plant hormone treatment has also shown effective results in regenerating wild apple populations. Wild apple trees were heavily lumbered for economical and agricultural uses in the mountains of Kazakhstan during the 1800s. Wild apple forests were turned into pastureland, which greatly changed the soil covering, and damaged young seedlings and roots. Prickly shrubs, such as eglantine and barberry exhibiting symbiotic relationships with wild apples by shielding them from predators were also cut. This further worsened the growing condition for wild apples, and severely weakened root suckers and therefore vegetative propagation. Cultivation These and other Malus species have been used in some recent breeding programmes to develop apples suitable for growing in harsh climates unsuitable for M. domestica, mainly for increased cold tolerance. A study in 2020 has discovered various gene inserts involved in dormancy and cold resistance features, such as heat shock proteins, in wild apples. In addition, desirable traits such as late flowering, early fruit maturity, short juvenility and stooling capability were studied by many breeding programs. Malus sieversii has recently been cultivated by the US Agricultural Research Service, in hopes of finding genetic information of value in the breeding of the modern apple plant. Some, but not all, of the resulting trees show unusual disease resistance. The variation in their response to disease on an individual basis is, itself, a sign of how much more genetically diverse they are than their domesticated descendants. For instance, wild apples were found to have multiple blue mold resistant genes, specifically against Penicillium expansum. The USDA Plant Genetic Resources Unit (PGRU) also conducted phenotypic analysis on M. sieversii seedlings, and has identified various pathogenic resistance including apple scab, fire blight, and cedar apple rust. A research in 2001 found various insect resistances within M. sieversii seedlings, and has identified instances for further research on its resistance for apple maggots and apple leaf curling midges. Effects of heat on M. sieversii were also studied in hot and arid regions, and they were found to be considerably drought tolerant and sunburn resistant. Red-fleshed apples Malus sieversii has been recently used as a critical source in the breeding of red-fleshed apples, due to its high genetic variability. This is seen as they are used to improve the stress resistance towards drought, cold, and pests of cultivated apple species. Some neglected characteristics of M. sieversii, such as high-flavonoid contents (especially anthocyanin) and short juvenile phases, have recently been used for red-fleshed apple breeding since traditional red-fleshed apples are not rich in these flavonoids. Using M. sieversii for breeding due to its high anthocyanin content has numerous benefits, including preventing cardiovascular disease and protecting against liver damage. The variant of M. sieversii, Malus sieversii f. niedzwetzkyana, has been emphasized for the use of breeding red-fleshed apples since it has red flowers, fruit skin, and pulp; in addition to its high anthocyanin content. When breeding Malus sieversii f. niedzwetzkyana, it was found that light results in higher anthocyanin production than those bred in the dark. Hybrids of Malus sieversii have also been an interest for breeders of red-fleshed apples. In culture Almaty, the largest city in Kazakhstan, and formerly its capital, derives its name from the Kazakh word alma 'apple', and it is often explained as meaning 'full of apples' (the region surrounding Almaty is home to forests of Malus sieversii). Gallery	Biology and health sciences	Pomes	Plants
3183582	https://en.wikipedia.org/wiki/Wire%20rope	Wire rope	Wire rope is composed of as few as two solid, metal wires twisted into a helix that forms a composite rope, in a pattern known as laid rope. Larger diameter wire rope consists of multiple strands of such laid rope in a pattern known as cable laid. Manufactured using an industrial machine known as a strander, the wires are fed through a series of barrels and spun into their final composite orientation. In stricter senses, the term wire rope refers to a diameter larger than , with smaller gauges designated cable or cords. Initially wrought iron wires were used, but today steel is the main material used for wire ropes. Historically, wire rope evolved from wrought iron chains, which had a record of mechanical failure. While flaws in chain links or solid steel bars can lead to catastrophic failure, flaws in the wires making up a steel cable are less critical as the other wires easily take up the load. While friction between the individual wires and strands causes wear over the life of the rope, it also helps to compensate for minor failures in the short run. Wire ropes were developed starting with mining hoist applications in the 1830s. Wire ropes are used dynamically for lifting and hoisting in cranes and elevators, and for transmission of mechanical power. Wire rope is also used to transmit force in mechanisms, such as a Bowden cable or the control surfaces of an airplane connected to levers and pedals in the cockpit. Only aircraft cables have WSC (wire strand core). Also, aircraft cables are available in smaller diameters than wire rope. For example, aircraft cables are available in diameter while most wire ropes begin at a diameter. Static wire ropes are used to support structures such as suspension bridges or as guy wires to support towers. An aerial tramway relies on wire rope to support and move cargo overhead. History Modern wire rope was invented by the German mining engineer Wilhelm Albert in the years between 1831 and 1834 for use in mining in the Harz Mountains in Clausthal, Lower Saxony, Germany. It was quickly accepted because it proved superior strength from ropes made of hemp or of metal chains, such as had been used before. Wilhelm Albert's first ropes consisted of three strands consisting of four wires each. In 1840, Scotsman Robert Stirling Newall improved the process further. In America wire rope was manufactured by John A. Roebling, starting in 1841 and forming the basis for his success in suspension bridge building. Roebling introduced a number of innovations in the design, materials and manufacture of wire rope. Ever with an ear to technology developments in mining and railroading, Josiah White and Erskine Hazard, principal owners of the Lehigh Coal & Navigation Company (LC&N Co.) — as they had with the first blast furnaces in the Lehigh Valley — built a Wire Rope factory in Jim Thorpe, Pennsylvania, in 1848, which provided lift cables for the Ashley Planes project, then the back track planes of the Summit Hill & Mauch Chunk Railroad, improving its attractiveness as a premier tourism destination, and vastly improving the throughput of the coal capacity since return of cars dropped from nearly four hours to less than 20 minutes. The following decades featured a burgeoning increase in deep shaft mining in both Europe and North America as surface mineral deposits were exhausted and miners had to chase layers along inclined layers. The era was early in railroad development and steam engines lacked sufficient tractive effort to climb steep slopes, so inclined plane railways were common. This pushed development of cable hoists rapidly in the United States as surface deposits in the Anthracite Coal Region north and south dove deeper every year, and even the rich deposits in the Panther Creek Valley required LC&N Co. to drive their first shafts into lower slopes beginning Lansford and its Schuylkill County twin-town Coaldale. The German engineering firm of Adolf Bleichert & Co. was founded in 1874 and began to build bicable aerial tramways for mining in the Ruhr Valley. With important patents, and dozens of working systems in Europe, Bleichert dominated the global industry, later licensing its designs and manufacturing techniques to Trenton Iron Works, New Jersey, USA which built systems across America. Adolf Bleichert & Co. went on to build hundreds of aerial tramways around the world: from Alaska to Argentina, Australia and Spitsbergen. The Bleichert company also built hundreds of aerial tramways for both the Imperial German Army and the Wehrmacht. In the latter part of the 19th century, wire rope systems were used as a means of transmitting mechanical power including for the new cable cars. Wire rope systems cost one-tenth as much and had lower friction losses than line shafts. Because of these advantages, wire rope systems were used to transmit power for a distance of a few miles or kilometers. Construction Wires Steel wires for wire ropes are normally made of non-alloy carbon steel with a carbon content of 0.4 to 0.95%. The very high strength of the rope wires enables wire ropes to support large tensile forces and to run over sheaves with relatively small diameters. Strands In the so-called cross lay strands, the wires of the different layers cross each other. In the mostly used parallel lay strands, the lay length of all the wire layers is equal and the wires of any two superimposed layers are parallel, resulting in linear contact. The wire of the outer layer is supported by two wires of the inner layer. These wires are neighbors along the whole length of the strand. Parallel lay strands are made in one operation. The endurance of wire ropes with this kind of strand is always much greater than of those (seldom used) with cross lay strands. Parallel lay strands with two wire layers have the construction Filler, Seale or Warrington. Spiral ropes In principle, spiral ropes are round strands as they have an assembly of layers of wires laid helically over a centre with at least one layer of wires being laid in the opposite direction to that of the outer layer. Spiral ropes can be dimensioned in such a way that they are non-rotating which means that under tension the rope torque is nearly zero. The open spiral rope consists only of round wires. The half-locked coil rope and the full-locked coil rope always have a centre made of round wires. The locked coil ropes have one or more outer layers of profile wires. They have the advantage that their construction prevents the penetration of dirt and water to a greater extent and it also protects them from loss of lubricant. In addition, they have one further very important advantage as the ends of a broken outer wire cannot leave the rope if it has the proper dimensions. Stranded ropes Stranded ropes are an assembly of several strands laid helically in one or more layers around a core. This core can be one of three types. The first is a fiber core, made up of synthetic material or natural fibers like sisal. Synthetic fibers are stronger and more uniform but cannot absorb much lubricant. Natural fibers can absorb up to 15% of their weight in lubricant and so protect the inner wires much better from corrosion than synthetic fibers do. Fiber cores are the most flexible and elastic, but have the downside of getting crushed easily. The second type, wire strand core, is made up of one additional strand of wire, and is typically used for suspension. The third type is independent wire rope core (IWRC), which is the most durable in all types of environments. Most types of stranded ropes only have one strand layer over the core (fibre core or steel core). The lay direction of the strands in the rope can be right (symbol Z) or left (symbol S) and the lay direction of the wires can be right (symbol z) or left (symbol s). This kind of rope is called ordinary lay rope if the lay direction of the wires in the outer strands is in the opposite direction to the lay of the outer strands themselves. If both the wires in the outer strands and the outer strands themselves have the same lay direction, the rope is called a lang lay rope (from Dutch contrary to , formerly Albert's lay or langs lay). Regular lay means the individual wires were wrapped around the centers in one direction and the strands were wrapped around the core in the opposite direction. Multi-strand ropes are all more or less resistant to rotation and have at least two layers of strands laid helically around a centre. The direction of the outer strands is opposite to that of the underlying strand layers. Ropes with three strand layers can be nearly non-rotating. Ropes with two strand layers are mostly only low-rotating. Classification according to usage Depending on where they are used, wire ropes have to fulfill different requirements. The main uses are: Running ropes (stranded ropes) are bent over sheaves and drums. They are therefore stressed mainly by bending and secondly by tension. Stationary ropes, stay ropes (spiral ropes, mostly full-locked) have to carry tensile forces and are therefore mainly loaded by static and fluctuating tensile stresses. Ropes used for suspension are often called cables. Track ropes (full locked ropes) have to act as rails for the rollers of cabins or other loads in aerial ropeways and cable cranes. In contrast to running ropes, track ropes do not take on the curvature of the rollers. Under the roller force, a so-called free bending radius of the rope occurs. This radius increases (and the bending stresses decrease) with the tensile force and decreases with the roller force. Wire rope slings (stranded ropes) are used to harness various kinds of goods. These slings are stressed by the tensile forces but first of all by bending stresses when bent over the more or less sharp edges of the goods. Rope drive Technical regulations apply to the design of rope drives for cranes, elevators, rope ways and mining installations. Factors that are considered in design include: Number of working cycles allowed before replacement or breakage of the rope Donandt force (yielding tensile force for a given bending diameter ratio /) - strict limit. The nominal rope tensile force must be smaller than the Donandt force . Rope safety factor, ratio between the rope's breaking strength and the maximum load to be expected Allowable number of broken strands before replacement Optimal rope diameter for a given sheave diameter, so as to obtain best working life The calculation of the rope drive limits depends on: Data of the used wire rope Rope tensile force Diameter of sheave or drum Simple bendings per working cycle Reverse bendings per working cycle Combined fluctuating tension and bending per working cycle Relative fluctuating tensile force Rope bending length Safety The wire ropes are stressed by fluctuating forces, by wear, by corrosion and in seldom cases by extreme forces. The rope life is finite and the safety is only ensured by inspection for the detection of wire breaks on a reference rope length, of cross-section loss, as well as other failures so that the wire rope can be replaced before a dangerous situation occurs. Installations should be designed to facilitate the inspection of the wire ropes. Lifting installations for passenger transportation require that a combination of several methods should be used to prevent a car from plunging downwards. Elevators must have redundant bearing ropes and a safety gear. Ropeways and mine hoistings must be permanently supervised by a responsible manager and the rope must be inspected by a magnetic method capable of detecting inner wire breaks. Terminations The end of a wire rope tends to fray readily, and cannot be easily connected to plant and equipment. There are different ways of securing the ends of wire ropes to prevent fraying. The common and useful type of end fitting for a wire rope is to turn the end back to form a loop. The loose end is then fixed back on the wire rope. Termination efficiencies vary from about 70% for a Flemish eye alone; to nearly 90% for a Flemish eye and splice; to 100% for potted ends and swagings. Thimbles When the wire rope is terminated with a loop, there is a risk that it will bend too tightly, especially when the loop is connected to a device that concentrates the load on a relatively small area. A thimble can be installed inside the loop to preserve the natural shape of the loop, and protect the cable from pinching and abrading on the inside of the loop. The use of thimbles in loops is industry best practice. The thimble prevents the load from coming into direct contact with the wires. Wire rope clips A wire rope clip, sometimes called a clamp, is used to fix the loose end of the loop back to the wire rope. It usually consists of a U-bolt, a forged saddle, and two nuts. The two layers of wire rope are placed in the U-bolt. The saddle is then fitted to the bolt over the ropes (the saddle includes two holes to fit to the U-bolt). The nuts secure the arrangement in place. Two or more clips are usually used to terminate a wire rope depending on the diameter. As many as eight may be needed for a diameter rope. The mnemonic "never saddle a dead horse" means that when installing clips, the saddle portion of the assembly is placed on the load-bearing or "live" side, not on the non-load-bearing or "dead" side of the cable. This is to protect the live or stress-bearing end of the rope against crushing and abuse. The flat bearing seat and extended prongs of the body are designed to protect the rope and are always placed against the live end. The US Navy and most regulatory bodies do not recommend the use of such clips as permanent terminations unless periodically checked and re-tightened. Eye splice or Flemish eye An eye splice may be used to terminate the loose end of a wire rope when forming a loop. The strands of the end of a wire rope are unwound a certain distance, then bent around so that the end of the unwrapped length forms an eye. The unwrapped strands are then plaited back into the wire rope, forming the loop, or an eye, called an eye splice. A Flemish eye, or Dutch Splice, involves unwrapping three strands (the strands need to be next to each other, not alternates) of the wire and keeping them off to one side. The remaining strands are bent around, until the end of the wire meets the "V" where the unwrapping finished, to form the eye. The strands kept to one side are now re-wrapped by wrapping from the end of the wire back to the "V" of the eye. These strands are effectively rewrapped along the wire in the opposite direction to their original lay. When this type of rope splice is used specifically on wire rope, it is called a "Molly Hogan", and, by some, a "Dutch" eye instead of a "Flemish" eye. Swaged terminations Swaging is a method of wire rope termination that refers to the installation technique. The purpose of swaging wire rope fittings is to connect two wire rope ends together, or to otherwise terminate one end of wire rope to something else. A mechanical or hydraulic swager is used to compress and deform the fitting, creating a permanent connection. Threaded studs, ferrules, sockets, and sleeves are examples of different swaged terminations. Swaging ropes with fibre cores is not recommended. Wedge sockets A wedge socket termination is useful when the fitting needs to be replaced frequently. For example, if the end of a wire rope is in a high-wear region, the rope may be periodically trimmed, requiring the termination hardware to be removed and reapplied. An example of this is on the ends of the drag ropes on a dragline. The end loop of the wire rope enters a tapered opening in the socket, wrapped around a separate component called the wedge. The arrangement is knocked in place, and load gradually eased onto the rope. As the load increases on the wire rope, the wedge become more secure, gripping the rope tighter. Potted ends or poured sockets Poured sockets are used to make a high strength, permanent termination; they are created by inserting the wire rope into the narrow end of a conical cavity which is oriented in-line with the intended direction of strain. The individual wires are splayed out inside the cone or 'capel', and the cone is then filled with molten lead–antimony–tin (Pb80Sb15Sn5) solder or 'white metal capping', zinc, or now more commonly, an unsaturated polyester resin compound.	Technology	Flexible components	null
3183595	https://en.wikipedia.org/wiki/Textile%20printing	Textile printing	Textile printing is the process of applying color to fabric in definite patterns or designs. In properly printed fabrics the colour is bonded with the fibre, so as to resist washing and friction. Textile printing is related to dyeing but in dyeing properly the whole fabric is uniformly covered with one colour, whereas in printing one or more colours are applied to it in certain parts only, and in sharply defined patterns. In printing, wooden blocks, stencils, engraved plates, rollers, or silkscreens can be used to place colours on the fabric. Colourants used in printing contain dyes thickened to prevent the colour from spreading by capillary attraction beyond the limits of a pattern or design. History Woodblock printing is a technique for printing text, images or patterns used widely throughout East Asia and probably originating in China in antiquity as a method of printing on textiles and later paper. As a method of printing on cloth, the earliest surviving examples from China date to before 220 CE/AD. Textile printing was known in Europe, via the Islamic world, from about the 12th century, and was widely used. However, the European dyes tended to liquify, which restricted the use of printed patterns. Fairly large and ambitious designs were printed for decorative purposes such as wall-hangings and lectern-cloths, where this was less of a problem as they did not need washing. When paper became common, the technology was rapidly used on that for woodcut prints. Superior cloth was also imported from Islamic countries, but this was much more expensive. Woodblock printing has a rich history and has evolved over time. In Japan, it became a highly respected art form called ukiyo-e, creating famous works like Hokusai's "The Great Wave Off Kanagawa" (circa 1830-1832) [1]. In Europe, it influenced other printing techniques like engraving and etching. Today, artists and craftspeople still use and innovate woodblock printing, keeping this ancient technique alive. The Incas of Peru, Chile and the Aztecs of Mexico also practiced textile printing previous to the Spanish Invasion in 1519. During the later half of the 17th century the French brought directly by sea, from their colonies on the east coast of India, samples of Indian blue and white resist prints, and along with them, particulars of the processes by which they had been produced, which produced washable fabrics. As early as the 1630s, the East India Company was bringing in printed and plain cotton for the English market. By the 1660s British printers and dyers were making their own printed cotton to sell at home, printing single colours on plain backgrounds; less colourful than the imported prints, but more to the taste of the British. Designs were also sent to India for their craftspeople to copy for export back to England. There were many dyehouses in England in the latter half of the 17th century, Lancaster being one area and on the River Lea near London another. Plain cloth was put through a prolonged bleaching process which prepared the material to receive and hold applied colour; this process vastly improved the colour durability of English calicoes and required a great deal of water from nearby rivers. One dyehouse was started by John Meakins, a London Quaker who lived in Cripplegate. When he died, he passed his dyehouse to his son-in-law Benjamin Ollive, Citizen and Dyer, who moved the dye-works to Bromley Hall where it remained in the family until 1823, known as Benjamin Ollive and Company, Ollive & Talwin, Joseph Talwin & Company and later Talwin & Foster. Samples of their fabrics and designs can be found in the Victoria and Albert Museum in London and the Smithsonian Copper-Hewett in New York. On the continent of Europe the commercial importance of calico printing seems to have been almost immediately recognized, and in consequence it spread and developed there much more rapidly than in England, where it was neglected for nearly ninety years after its introduction. During the last two decades of the 17th century and the earlier ones of the 18th new dye works were started in France, Germany, Switzerland and Austria. It was only in 1738 that calico printing was first, practiced in Scotland, and not until twenty-six years later that Messrs Clayton of Bamber Bridge, near Preston, established in 1764 the first print-works in Lancashire, and thus laid the foundation of the industry. From an artistic point of view most of the pioneer work in calico printing was done by the French. From the early days of the industry down to the latter half of the 20th century, the productions of the French printers in Jouy, Beauvais, Rouen, and in Alsace-Lorraine, were looked upon as representing all that was best in artistic calico printing. Methods Traditional textile printing techniques may be broadly categorized into four styles: Direct printing, in which colourants containing dyes, thickeners, and the mordants or substances necessary for fixing the colour on the cloth are printed in the desired pattern. The printing of a mordant in the desired pattern prior to dyeing cloth; the colour adheres only where the mordant was printed. Resist dyeing, in which a wax or other substance is printed onto fabric which is subsequently dyed. The waxed areas do not accept the dye, leaving uncoloured patterns against a coloured ground. Discharge printing, in which a bleaching agent is printed onto previously dyed fabrics to remove some or all of the colour. Resist and discharge techniques were particularly fashionable in the 19th century, as were combination techniques in which indigo resist was used to create blue backgrounds prior to block-printing of other colours. Modern industrial printing mainly uses direct printing techniques. The printing process does involve several stages in order to prepare the fabric and printing paste, and to fix the impression permanently on the fabric: pre-treatment of fabric, preparation of colours, preparation of printing paste, impression of paste on fabric using printing methods, drying of fabric, fixing the printing with steam or hot air (for pigments), after process treatments. Preparation of cloth for printing Cloth is prepared by washing and bleaching. For a coloured ground it is then dyed. The cloth has always to be brushed, to free it from loose nap, flocks and dust that it picks up whilst stored. Frequently, too, it has to be sheared by being passed over rapidly revolving knives arranged spirally round an axle, which rapidly and effectually cuts off all filaments and knots, leaving the cloth perfectly smooth and clean and in a condition fit to receive impressions of the most delicate engraving. Some fabrics require very careful stretching and straightening on a stenter before they are wound around hollow wooden or iron centers into rolls of convenient size for mounting on the printing machines. Preparation of colours The art of making colours for textile printing demands both chemical knowledge and extensive technical experience, for their ingredients must not only be in proper proportion to each other, but also specially chosen and compounded for the particular style of work in hand. A colour must comply to conditions such as shade, quality and fastness; where more colours are associated in the same design each must be capable of withstanding the various operations necessary for the development and fixation of the others. All printing pastes whether containing colouring matter or not are known technically as colours. Colours vary considerably in composition. Most of them contain all the elements necessary for direct production and fixation. Some, however, contain the colouring matter alone and require various after-treatments; and others again are simply thickened mordants. A mordant is a metallic salt or other substance that combines with the dye to form an insoluble colour, either directly by steaming, or indirectly by dyeing. All printing colours require thickening to enable them to be transferred from colour-box to cloth without running or spreading beyond the limits of the pattern. Thickening agents The printing thickeners used depend on the printing technique, the fabric and the particular dyestuff . Typical thickening agents are starch derivatives, flour, gum arabic, guar gum derivatives, tamarind, sodium alginate, sodium polyacrylate, gum Senegal and gum tragacanth, British gum or dextrin and albumen. Hot-water-soluble thickening agents such as native starch are made into pastes by boiling in double or jacketed pans. Most thickening agents used today are cold-soluble and require only extensive stirring. Starch paste Starch paste is made from wheat starch, cold water, and olive oil, then thickened by boiling. Non-modified starch is applicable to all but strongly alkaline or strongly acid colours. With the former it thickens up to a stiff unworkable jelly. In the case of the latter, while mineral acids or acid salts convert it into dextrin, thus diminishing its viscosity or thickening power, organic acids do not have that effect. Today, modified carboxymethylated cold soluble starches are mainly used. These have a stable viscosity and are easy to rinse out of the fabric and give reproducible "short" paste rheology. Flour paste is made in a similar way to starch paste; it is sometimes used to thicken aluminum and iron mordants. Starch paste resists of rice flour have been used for several centuries in Japan. Gums Gum arabic and gum Senegal are both traditional thickenings, but expense prevents them from being used for any but pale, delicate tints. They are especially useful thickenings for the light ground colours of soft muslins and 9 penetrate as well into the fibre of the cloth or as deeply as pure starch or flour and is unsuitable for very dark, strong colours. Gum tragacanth, or Dragon, which may be mixed in any proportion with starch or flour, is equally useful for pigment colours and mordant colours. When added to a starch paste it increases its penetrative power and adds to its softness without diminishing its thickness, making it easier to wash out of the fabric. It produces much more even colours than does starch paste alone. Used by itself it is suitable for printing all kinds of dark grounds on goods that are required to retain their soft "clothy" feel. Starch always leaves the printed cloth somewhat harsh in feeling (unless modified carboxymethylated starches are used), but very dark colours can be obtained. Gum Senegal, gum arabic or modified guar gum thickening yield clearer and more even tints than does starch, suitable for lighter colours but less suited for very dark colours. (The gums apparently prevent the colours from combining fully with the fibers.) A printing stock solution is mostly a combination of modified starch and gum stock solutions. Albumen Albumen is both a thickening and a fixing agent for insoluble pigments. Chrome yellow, the ochres, vermilion and ultramarine are such pigments. Albumen is always dissolved in the cold, a process that takes several days when large quantities are required. Egg albumen is expensive and only used for the lightest shades. Blood albumen solution is used in cases when very dark colours are required to be absolutely fast to washing. After printing, albumen thickened colours are exposed to hot steam, which coagulates the albumen and effectually fixes the colours. Printing paste preparation Combinations of cold water-soluble carboxymethylated starch, guar gum and tamarind derivatives are most commonly used today in disperse screen printing on polyester. Alginates are used for cotton printing with reactive dyes, sodium polyacrylates for pigment printing, and in the case of vat dyes on cotton only carboxymethylated starch is used. Formerly, colours were always prepared for printing by boiling the thickening agent, the colouring matter and solvents, together, then cooling and adding various fixing agents. At the present time, however, concentrated solutions of the colouring matters and other adjuncts are often simply added to the cold thickenings, of which large quantities are kept in stock. Colours are reduced in shade by simply adding more stock (printing) paste. For example, a dark blue containing 4 oz. of methylene blue per gallon may readily be made into a pale shade by adding to it thirty times its bulk of starch paste or gum, as the case may be. The procedure is similar for other colours. Before printing it is essential to strain or sieve all colours in order to free them from lumps, fine sand, and other impurities, which would inevitably damage the highly polished surface of the engraved rollers and result in bad printing. Every scratch on the surface of a roller prints a fine line on the cloth, and too much care, therefore, cannot be taken to remove, as far as possible, all grit and other hard particles from every colour. Straining is usually done by squeezing the colour through filter cloths like artisanal fine cotton, silk or industrial woven nylon. Fine sieves can also be employed for colours that are used hot or are very strongly alkaline or acid. Methods of printing There are eight distinct methods presently used to impress coloured patterns on cloth: Hand block printing Perrotine printing Engraved copperplate printing Roller, cylinder, or machine printing Stencil printing Screen printing Digital textile printing Flexo textile printing Discharge Printing Heat transfer printing Gallery Block printing This process is the earliest, simplest and slowest of all printing methods. A design is drawn on, or transferred to, prepared wooden blocks. A separate block is required for each distinct colour in the design. A block cutter carves out the wood around the heavier masses first, leaving the finer and more delicate work until the last so as to avoid any risk of injuring it when the coarser parts are cut. When finished, the block has the appearance of a flat relief carving, with the design standing out. Fine details, difficult to cut in wood, are built up in strips of brass or copper, which is bent to shape and driven edgewise into the flat surface of the block. This method is known as coppering. The printer applies colour to the block and presses it firmly and steadily on the cloth, striking it smartly on the back with a wooden mallet. The second impression is made in the same way, the printer taking care to see that it registers exactly with the first. Pins at each corner of the block join up exactly, so that the pattern can continue without a break. Each succeeding impression is made in precisely the same manner until the length of cloth is fully printed. The cloth is then wound over drying rollers. If the pattern contains several colours the cloth is first printed throughout with one colour, dried, and then printed with the next. Block printing by hand is a slow process. It is, however, capable of yielding highly artistic results, some of which are unobtainable by any other method. William Morris used this technique in some of his fabrics. Perrotine printing The perrotine is a block-printing machine invented by Perrot of Rouen in 1834 and is now only of historical interest. Roller, cylinder, or machine printing This process was patented by Thomas Bell in 1785, fifteen years after his use of an engraved plate to print textiles. Bell's patent was for a machine to print six colours at once, but, probably owing to its incomplete development, it was not immediately successful. One colour could be printed with satisfactorily; the difficulty was to keep the six rollers in register with each other. This defect was overcome by Adam Parkinson of Manchester in 1785. That year, Bells machine with Parkinson's improvement was successfully employed by Messrs Livesey, Hargreaves and Company of Bamber Bridge, Preston, for the printing of calico in from two to six colours at a single operation. Roller printing was highly productive, 10,000 to 12,000 yards being commonly printed in one day of ten hours by a single-colour machine. It is capable of reproducing every style of design, ranging from the fine delicate lines of copperplate engraving to the small repeats and limited colours of the perrotine to the broadest effects of block printing with repeats from 1 in to 80 inches. It is precise, so each portion of an elaborate multicolour pattern can be fitted into its proper place without faulty joints at the points of repetition. Stencil printing The art of stenciling on textile fabrics has been practiced from time immemorial by the Japanese, and found increasing employment in Europe for certain classes of decorative work on woven goods during the late 19th century. A pattern is cut from a sheet of stout paper or thin metal with a sharp-pointed knife, the uncut portions representing the part that will be left uncoloured. The sheet is laid on the fabric and colour is brushed through its interstices. The peculiarity of stenciled patterns is that they have to be held together by ties. For instance, a complete circle cannot be cut without its centre dropping out, so its outline has to be interrupted at convenient points by ties or uncut portions. This limitation influences the design. For single-colour work a stenciling machine was patented in 1894 by S. H. Sharp. It consists of an endless stencil plate of thin sheet steel that passes continuously over a revolving cast iron cylinder. The cloth to be ornamented passes between the two and the colour is forced onto it through the holes in the stencil by mechanical means. Screen-printing Screen printing is by far the most common technology today. Two types exist: rotary screen printing and flat (bed) screen printing. A blade (squeegee) squeezes the printing paste through openings in the screen onto the fabric. Digital textile printing Digital textile printing is often referred to as direct-to-garment printing (DTG printing), or digital garment printing. It is a process of printing on textiles and garments using specialized or modified inkjet technology. Inkjet printing on fabric is also possible with an inkjet printer by using fabric sheets with a removable paper backing. Today, major inkjet technology manufacturers can offer specialized products designed for direct printing on textiles, not only for sampling but also for bulk production. Since the early 1990s, inkjet technology and specially developed water-based ink (known as dye-sublimation or disperse direct ink) have made it possible to print directly onto polyester fabric. This is mainly related to visual communication in retail and brand promotion (flags, banners and other point of sales applications). Printing onto nylon and silk can be done by using an acid ink. Reactive ink is used for cellulose based fibers such as cotton and linen. Inkjet technology in digital textile printing allows for single pieces, mid-run production and even long-run alternatives to screen printed fabric. A similar printing method: Direct-To-Film printing (DTF printing) can also make the digital textile printing. The difference from DTG printing is that DTF printing first prints on a special transfer film while DTG printing prints on the substrate. One of the advantage of DTF printing is that it is more cost effective. Flexo textile printing Flexo textile printing on textile fabric was successful in China in the last 4 years. Central Impression Flexo, Rubber Sleeves as the printing plate in round engraved by laser (Direct Laser Engraving), Anilox in Sleeve technologies are applicated in the area. Not only the solid, but also 6 to 8 colours in fine register, higher resolution ratio and higher productivity which are the outstanding advantages extraordinary different from the traditional screen textile printing. Aerospace Huayang, Hell system, SPGPrints and Felix Böttcher contributed their technologies and efforts. Other methods of printing Although most work is executed throughout by one or another of the seven distinct processes mentioned above, combinations are frequently employed. Sometimes a pattern is printed partly by machine and partly by block, and sometimes a cylindrical block is used along with engraved copper-rollers in an ordinary printing machine. The block in this latter case is in all respects, except for shape, identical with a flat wood or coppered block, but, instead of being dipped in colour, it receives its supply from an endless blanket, one part of which works in contact with colour-furnishing rollers and the other part with the cylindrical block. This block is known as a surface or peg roller. Many attempts have been made to print multicolour patterns with surface rollers alone, but hitherto with little success, owing to their irregularity in action and to the difficulty of preventing them from warping. These defects are not present in the printing of linoleum in which opaque oil colours are used, colours that neither sink into the body of the hard linoleum nor tend to warp the roller. Lithographic printing has been applied to textile fabrics with qualified success. Its irregularity and the difficulty of registering repeats have restricted its use to the production of decorative panels, equal or smaller in size to the plate or stone. Calico printing Goods intended for calico printing are well-bleached; otherwise stains and other serious defects are certain to arise during subsequent operations. The chemical preparations used for special styles will be mentioned in their proper places; but a general prepare, employed for most colours that are developed and fixed by steaming only, consists in passing the bleached calico through a weak solution of sulphated or turkey red oil containing 2.5 to 5 percent fatty acid. Some colours are printed on pure bleached cloth, but all patterns containing alizarin red, rose and salmon shades are considerably brightened by the presence of oil, and indeed very few, if any, colours are detrimentally affected by it. The cloth is always brushed to free it from loose nap, flocks and dust that it picks up whilst stored. Frequently, too, it has to be sheared by being passed over rapidly revolving knives arranged spirally round an axle, which rapidly and effectually cuts off all filaments and knots, leaving the cloth perfectly smooth and clean. It is then stentered, wound onto a beam, and mounting on the printing machines. Woollen printing The printing of wool differs little from the printing of cotton in general. Most of the colours employed in the one industry are used in the other, and the operations of steaming, washing and soaping are almost identical. Unlike cotton, however, wool requires special preparation, after bleaching, if the full tinctorial value of the colours is to be obtained. Two quite different methods of preparation are used, namely (1) the chlorination of the wool; and (2) the precipitation of stannic acid on the fibre. In the first method the woollen fabric is first passed through a solution of bleaching powder, then well squeezed and passed, without washing, into dilute sulphuric or hydrochloric acid, squeezed again and well washed in water, after which it is dried. Great care and experience are demanded in this operation to prevent the wool from becoming hard and yellow. In the second method the cloth is padded in sodium stannate, well squeezed, passed into dilute sulphuric acid, well washed and dried. For certain styles of work it is necessary to combine both preparations. Although alizarin, mordant colours and dyewood extracts can be used on wool, the vast majority of patterns printed on wool are executed by means of acid dye-stuffs and basic colours, for both of which this fibre possesses a natural affinity. In most cases therefore these colours are simply dissolved in a little acetic and citric acids, thickened with gum and printed without any further addition. The addition of tannic acid, however, can be made to, and considerably increases the fastness of, the basic dyes. Mordant colours like logwood black are applied in the usual way. The printing of wool is carried out exactly as for cotton, but if the best results are to be obtained, the engraving of the rollers must be deep, the blanket on the machine as solt as possible, and the drying of the printed cloth very gentle. After printing, the goods are steamed in moist steam or wrapped between moistened "greys" and steamed in a "cottage" steamer. If too little moisture is given, the colours lack both strength and brilliancy; if too much they run. The correct degree of dampness can only he determined by experience of the work, combined with a special knowledge of the particular apparatus employed. After steaming, the printed goods are washed in plenty of water, then dried up and. finished with a little glycerol or some waxy preparation. Discharges may be very easily obtained on wool dyed in acid dye-stuffs, by means of stannous chloride and basic colours for the coloured effect, and hydrosulphite for the white. Silk printing The colours and methods employed are the same as for wool, except that in the case of silk no preparation of the material is required before printing, and ordinary dry steaming is preferable to damp steaming. Both acid and basic dyes play an important role in silk printing, which for the most part is confined to the production of articles for fashion goods, handkerchiefs, and scarves, all articles for which bright colours are in demand. Alizarine and other mordant colours are mainly used for any goods that have to resist repeated washings or prolonged exposure to light. In this case the silk frequently must be prepared in alizarine oil, after which it is treated in all respects like cotton, namely steamed, washed and soaped, the colours used being the same. Silk is especially adapted to discharge and reserve effects. Most of the acid dyes can be discharged in the same way as when they are dyed on wool. Reserved effects are produced by printing mechanical resists, such as waxes and fats, on the cloth and then dyeing it in cold dye-liquor. The great affinity of the silk fibre for basic and acid dyestuffs enables it to extract colouring matter from cold solutions and permanently combine with it to form an insoluble lake. After dyeing, the reserve prints are washed, first in cold water to remove any colour not fixed onto the fibre, and then in hot water or benzene to dissolve out the resisting bodies. After steaming, silk goods are normally only washed in hot water, but those printed entirely in mordant dyes will stand soaping, and indeed require it to brighten the colours and soften the material. Some silk dyes do not require heat setting or steaming. They strike instantly, allowing the designer to dye colour upon colour. These dyes are intended mostly for silk scarf dyeing. They also dye bamboo, rayon, linen, and some other natural fabrics like hemp and wool to a lesser extent, but do not set on cotton. Artificial fibre printing	Technology	Other techniques	null
3184580	https://en.wikipedia.org/wiki/Solid-state%20laser	Solid-state laser	A solid-state laser is a laser that uses a gain medium that is a solid, rather than a liquid as in dye lasers or a gas as in gas lasers. Semiconductor-based lasers are also in the solid state, but are generally considered as a separate class from solid-state lasers, called laser diodes. Solid-state media Generally, the active medium of a solid-state laser consists of a glass or crystalline "host" material, to which is added a "dopant" such as neodymium, chromium, erbium, thulium or ytterbium. Many of the common dopants are rare-earth elements, because the excited states of such ions are not strongly coupled with the thermal vibrations of their crystal lattices (phonons), and their operational thresholds can be reached at relatively low intensities of laser pumping. There are many hundreds of solid-state media in which laser action has been achieved, but relatively few types are in widespread use. Of these, probably the most common is neodymium-doped yttrium aluminum garnet (Nd:YAG). Neodymium-doped glass (Nd:glass) and ytterbium-doped glasses or ceramics are used at very high power levels (terawatts) and high energies (megajoules), for multiple-beam inertial confinement fusion. The first material used for lasers was synthetic ruby crystals. Ruby lasers are still used for a few applications, but they are no longer common because of their low power efficiencies. At room temperature, ruby lasers emit only short pulses of light, but at cryogenic temperatures they can be made to emit a continuous train of pulses. The second solid-state gain medium was uranium-doped calcium fluoride. Peter Sorokin and Mirek Stevenson at IBM's laboratories in Yorktown Heights (US) experimented with this material in the 1960s and achieved lasing at 2.5 μm shortly after Maiman's ruby laser. Some solid-state lasers can be made tunable by using intracavity etalons, prisms, gratings, or a combination of these. Titanium-doped sapphire is widely used for its broad tuning range, 660 to 1080 nanometers. Alexandrite lasers are tunable from 700 to 820 nm and yield higher-energy pulses than titanium-sapphire lasers because of the gain medium's longer energy storage time and higher damage threshold. Pumping Solid state lasing media are typically optically pumped, using either a flashlamp or arc lamp, or by laser diodes. Diode-pumped solid-state lasers tend to be much more efficient and have become much more common as the cost of high-power semiconductor lasers has decreased. Mode locking Mode locking of solid-state lasers and fiber lasers has wide applications as large-energy ultra-short pulses can be obtained. There are two types of saturable absorbers that are widely used as mode lockers: SESAM, and SWCNT. Graphene has also been used. These materials use a nonlinear optical behavior called saturable absorption to make a laser create short pulses. Applications Solid state lasers are used in research, medical treatment, and military applications, among others.	Technology	Lasers	null
2307285	https://en.wikipedia.org/wiki/NASA/IPAC%20Extragalactic%20Database	NASA/IPAC Extragalactic Database	The NASA/IPAC Extragalactic Database (NED) is an online astronomical database for astronomers that collates and cross-correlates astronomical information on extragalactic objects (galaxies, quasars, radio, x-ray and infrared sources, etc.). NED was created in the late 1980s by two Pasadena astronomers, George Helou and Barry F. Madore. NED is funded by NASA and is operated by the Infrared Processing and Analysis Center (IPAC) on the campus of the California Institute of Technology, under contract with NASA. NED is built around a master list of extragalactic objects for which cross-identifications of names have been established, accurate positions and redshifts entered to the extent possible, and some basic data collected. Bibliographic references relevant to individual objects have been compiled, and abstracts of extragalactic interest are kept on line. Detailed and referenced photometry, position, and redshift data, have been taken from large compilations and from the literature. NED also includes images from 2MASS, from the literature, and from the Digitized Sky Survey. As of March 2014, NED contains 206 million distinct astronomical objects with 232 million cross-identifications across multiple wavelengths, with redshift measurements for 5 million objects, 1.9 billion photometric data points, 609 million diameter measurements, 71 thousand redshift-independent distances for over 15 thousand galaxies, 310 thousand detailed classifications for 230 thousand objects, and 2.6 million images, maps and external links, together with links to 65 thousand journal articles, notes and abstracts.	Physical sciences	Databases	Astronomy
2307295	https://en.wikipedia.org/wiki/SIMBAD	SIMBAD	SIMBAD (the Set of Identifications, Measurements and Bibliography for Astronomical Data) is an astronomical database of objects beyond the Solar System. It is maintained by the Centre de données astronomiques de Strasbourg (CDS), France. SIMBAD was created by merging the Catalog of Stellar Identifications (CSI) and the Bibliographic Star Index as they existed at the Meudon Computer Centre until 1979, and then expanded by additional source data from other catalogues and the academic literature. The first on-line interactive version, known as Version 2, was made available in 1981. Version 3, developed in the C language and running on UNIX stations at the Volgograd Observatory, was released in 1990. Fall of 2006 saw the release of Version 4 of the database, now stored in PostgreSQL, and the supporting software, now written entirely in Java. , SIMBAD contains information for 5,800,000 stars and about 5,500,000 nonstellar objects (galaxies, planetary nebulae, clusters, novae and supernovae, etc.). The minor planet 4692 SIMBAD was named in its honour.	Physical sciences	Databases	Astronomy
2307854	https://en.wikipedia.org/wiki/Rotations%20in%204-dimensional%20Euclidean%20space	Rotations in 4-dimensional Euclidean space	In mathematics, the group of rotations about a fixed point in four-dimensional Euclidean space is denoted SO(4). The name comes from the fact that it is the special orthogonal group of order 4. In this article rotation means rotational displacement. For the sake of uniqueness, rotation angles are assumed to be in the segment except where mentioned or clearly implied by the context otherwise. A "fixed plane" is a plane for which every vector in the plane is unchanged after the rotation. An "invariant plane" is a plane for which every vector in the plane, although it may be affected by the rotation, remains in the plane after the rotation. Geometry of 4D rotations Four-dimensional rotations are of two types: simple rotations and double rotations. Simple rotations A simple rotation about a rotation centre leaves an entire plane through (axis-plane) fixed. Every plane that is completely orthogonal to intersects in a certain point . For each such point is the centre of the 2D rotation induced by in . All these 2D rotations have the same rotation angle . Half-lines from in the axis-plane are not displaced; half-lines from orthogonal to are displaced through ; all other half-lines are displaced through an angle less than . Double rotations For each rotation of 4-space (fixing the origin), there is at least one pair of orthogonal 2-planes and each of which is invariant and whose direct sum is all of 4-space. Hence operating on either of these planes produces an ordinary rotation of that plane. For almost all (all of the 6-dimensional set of rotations except for a 3-dimensional subset), the rotation angles in plane and in plane – both assumed to be nonzero – are different. The unequal rotation angles and satisfying , are almost uniquely determined by . Assuming that 4-space is oriented, then the orientations of the 2-planes and can be chosen consistent with this orientation in two ways. If the rotation angles are unequal (), is sometimes termed a "double rotation". In that case of a double rotation, and are the only pair of invariant planes, and half-lines from the origin in , are displaced through and respectively, and half-lines from the origin not in or are displaced through angles strictly between and . Isoclinic rotations If the rotation angles of a double rotation are equal then there are infinitely many invariant planes instead of just two, and all half-lines from are displaced through the same angle. Such rotations are called isoclinic or equiangular rotations, or Clifford displacements. Beware: not all planes through are invariant under isoclinic rotations; only planes that are spanned by a half-line and the corresponding displaced half-lines are invariant. Assuming that a fixed orientation has been chosen for 4-dimensional space, isoclinic 4D rotations may be put into two categories. To see this, consider an isoclinic rotation , and take an orientation-consistent ordered set of mutually perpendicular half-lines at (denoted as ) such that and span an invariant plane, and therefore and also span an invariant plane. Now assume that only the rotation angle is specified. Then there are in general four isoclinic rotations in planes and with rotation angle , depending on the rotation senses in and . We make the convention that the rotation senses from to and from to are reckoned positive. Then we have the four rotations , , and . and are each other's inverses; so are and . As long as lies between 0 and , these four rotations will be distinct. Isoclinic rotations with like signs are denoted as left-isoclinic; those with opposite signs as right-isoclinic. Left- and right-isoclinic rotations are represented respectively by left- and right-multiplication by unit quaternions; see the paragraph "Relation to quaternions" below. The four rotations are pairwise different except if or . The angle corresponds to the identity rotation; corresponds to the central inversion, given by the negative of the identity matrix. These two elements of SO(4) are the only ones that are simultaneously left- and right-isoclinic. Left- and right-isocliny defined as above seem to depend on which specific isoclinic rotation was selected. However, when another isoclinic rotation with its own axes , , , is selected, then one can always choose the order of , , , such that can be transformed into by a rotation rather than by a rotation-reflection (that is, so that the ordered basis , , , is also consistent with the same fixed choice of orientation as , , , ). Therefore, once one has selected an orientation (that is, a system of axes that is universally denoted as right-handed), one can determine the left or right character of a specific isoclinic rotation. Group structure of SO(4) SO(4) is a noncommutative compact 6-dimensional Lie group. Each plane through the rotation centre is the axis-plane of a commutative subgroup isomorphic to SO(2). All these subgroups are mutually conjugate in SO(4). Each pair of completely orthogonal planes through is the pair of invariant planes of a commutative subgroup of SO(4) isomorphic to . These groups are maximal tori of SO(4), which are all mutually conjugate in SO(4).	Mathematics	Four-dimensional space	null
2309427	https://en.wikipedia.org/wiki/Magnetosphere%20of%20Saturn	Magnetosphere of Saturn	The magnetosphere of Saturn is the cavity created in the flow of the solar wind by the planet's internally generated magnetic field. Discovered in 1979 by the Pioneer 11 spacecraft, Saturn's magnetosphere is the second largest of any planet in the Solar System after Jupiter. The magnetopause, the boundary between Saturn's magnetosphere and the solar wind, is located at a distance of about 20 Saturn radii from the planet's center, while its magnetotail stretches hundreds of Saturn radii behind it. Saturn's magnetosphere is filled with plasmas originating from both the planet and its moons. The main source is the small moon Enceladus, which ejects as much as 1,000 kg/s of water vapor from the geysers on its south pole, a portion of which is ionized and forced to co-rotate with the Saturn's magnetic field. This loads the field with as much as 100 kg of water group ions per second. This plasma gradually moves out from the inner magnetosphere via the interchange instability mechanism and then escapes through the magnetotail. The interaction between Saturn's magnetosphere and the solar wind generates bright oval aurorae around the planet's poles observed in visible, infrared and ultraviolet light. The aurorae are related to the powerful saturnian kilometric radiation (SKR), which spans the frequency interval between 100 kHz to 1300 kHz and was once thought to modulate with a period equal to the planet's rotation. However, later measurements showed that the periodicity of the SKR's modulation varies by as much as 1%, and so probably does not exactly coincide with Saturn's true rotational period, which as of 2010 remains unknown. Inside the magnetosphere there are radiation belts, which house particles with energy as high as tens of megaelectronvolts. The energetic particles have significant influence on the surfaces of inner icy moons of Saturn. In 1980–1981 the magnetosphere of Saturn was studied by the Voyager spacecraft. Up until September 2017 it was a subject of ongoing investigation by Cassini mission, which arrived in 2004 and spent over 13 years observing the planet. Discovery Immediately after the discovery of Jupiter's decametric radio emissions in 1955, attempts were made to detect a similar emission from Saturn, but with inconclusive results. The first evidence that Saturn might have an internally generated magnetic field came in 1974, with the detection of weak radio emissions from the planet at the frequency of about 1 MHz. These medium wave emissions were modulated with a period of about , which was interpreted as Saturn's rotation period. Nevertheless, the evidence available in the 1970s was too inconclusive and some scientists thought that Saturn might lack a magnetic field altogether, while others even speculated that the planet could lie beyond the heliopause. The first definite detection of the saturnian magnetic field was made only on September 1, 1979, when it was passed through by the Pioneer 11 spacecraft, which measured its magnetic field strength directly. Structure Internal field Like Jupiter's magnetic field, Saturn's is created by a fluid dynamo within a layer of circulating liquid metallic hydrogen in its outer core. Like Earth, Saturn's magnetic field is mostly a dipole, with north and south poles at the ends of a single magnetic axis. On Saturn, like on Jupiter, the north magnetic pole is located in the northern hemisphere, and the south magnetic pole lies in the southern hemisphere, so that magnetic field lines point away from the north pole and towards the south pole. This is reversed compared to the Earth, where the north magnetic pole lies in the southern hemisphere. Saturn's magnetic field also has quadrupole, octupole and higher components, though they are much weaker than the dipole. The magnetic field strength at Saturn's equator is about 21 μT (0.21 G), which corresponds to a dipole magnetic moment of about 4.6 T•m3. This makes Saturn's magnetic field slightly weaker than Earth's; however, its magnetic moment is about 580 times larger. Saturn's magnetic dipole is strictly aligned with its rotational axis, meaning that the field, uniquely, is highly axisymmetric. The dipole is slightly shifted (by 0.037 Rs) along Saturn's rotational axis towards the north pole. Size and shape Saturn's internal magnetic field deflects the solar wind, a stream of ionized particles emitted by the Sun, away from its surface, preventing it from interacting directly with its atmosphere and instead creating its own region, called a magnetosphere, composed of a plasma very different from that of the solar wind. The magnetosphere of Saturn is the second–largest magnetosphere in the Solar System after that of Jupiter. As with Earth's magnetosphere, the boundary separating the solar wind's plasma from that within Saturn's magnetosphere is called the magnetopause. The magnetopause distance from the planet's center at the subsolar point varies widely from 16 to 27 Rs (Rs=60,330 km is the equatorial radius of Saturn). The magnetopause's position depends on the pressure exerted by the solar wind, which in turn depends on solar activity. The average magnetopause standoff distance is about 22 Rs. In front of the magnetopause (at the distance of about 27 Rs from the planet) lies the bow shock, a wake-like disturbance in the solar wind caused by its collision with the magnetosphere. The region between the bow shock and magnetopause is called the magnetosheath. At the opposite side of the planet, the solar wind stretches Saturn's magnetic field lines into a long, trailing magnetotail, which consists of two lobes, with the magnetic field in the northern lobe pointing away from Saturn and the southern pointing towards it. The lobes are separated by a thin layer of plasma called the tail current sheet. Like Earth's, Saturn's tail is a channel through which solar plasma enters the inner regions of the magnetosphere. Similar to Jupiter, the tail is the conduit through which the plasma of the internal magnetospheric origin leaves the magnetosphere. The plasma moving from the tail to the inner magnetosphere is heated and forms a number of radiation belts. Magnetospheric regions Saturn's magnetosphere is often divided into four regions. The innermost region co-located with Saturn's planetary rings, inside approximately 3 Rs, has a strictly dipolar magnetic field. It is largely devoid of plasma, which is absorbed by ring particles, although the radiation belts of Saturn are located in this innermost region just inside and outside the rings. The second region between 3 and 6 Rs contains the cold plasma torus and is called the inner magnetosphere. It contains the densest plasma in the saturnian system. The plasma in the torus originates from the inner icy moons and particularly from Enceladus. The magnetic field in this region is also mostly dipolar. The third region lies between 6 and 12–14 Rs and is called the dynamic and extended plasma sheet. The magnetic field in this region is stretched and non-dipolar, whereas the plasma is confined to a thin equatorial plasma sheet. The fourth outermost region is located beyond 15 Rs at high latitudes and continues up to magnetopause boundary. It is characterized by a low plasma density and a variable, non-dipolar magnetic field strongly influenced by the Solar wind. In the outer parts of Saturn's magnetosphere beyond approximately 15–20 Rs the magnetic field near the equatorial plane is highly stretched and forms a disk-like structure called magnetodisk. The disk continues up to the magnetopause on the dayside and transitions into the magnetotail on the nightside. Near the dayside it can be absent when the magnetosphere is compressed by the Solar wind, which usually happens when the magnetopause distance is smaller than 23 Rs. On the nightside and flanks of the magnetosphere the magnetodisk is always present. The Saturnian magnetodisk is a much smaller analog of the Jovian magnetodisk. The plasma sheet in Saturn's magnetosphere has a bowl-like shape not found in any other known magnetosphere. When Cassini arrived in 2004, there was a winter in the northern hemisphere. The measurements of the magnetic field and plasma density revealed that the plasma sheet was warped and lay to the north of the equatorial plane, looking like a giant bowl. Such a shape was unexpected. Dynamics The processes driving Saturn's magnetosphere are similar to those driving Earth's and Jupiter's. Just as Jupiter's magnetosphere is dominated by plasma co–rotation and mass–loading from Io, so Saturn's magnetosphere is dominated by plasma co–rotation and mass–loading from Enceladus. However, Saturn's magnetosphere is much smaller in size, while its inner region contains too little plasma to seriously distend it and create a large magnetodisk. This means that it is much more strongly influenced by the solar wind, and that, like Earth's magnetic field, its dynamics are affected by reconnection with the wind similar to the Dungey cycle. Another distinguishing feature of Saturn's magnetosphere is high abundance of neutral gas around the planet. As revealed by ultraviolet observation of Cassini, the planet is enshrouded in a large cloud of hydrogen, water vapor and their dissociative products like hydroxyl, extending as far as 45 Rs from Saturn. In the inner magnetosphere the ratio of neutrals to ions is around 60 and it increases in the outer magnetosphere, which means that the entire magnetospheric volume is filled with relatively dense weakly ionized gas. This is different, for instance, from Jupiter or Earth, where ions dominate over neutral gas, and has consequences for the magnetospheric dynamics. Sources and transport of plasma The plasma composition in Saturn's inner magnetosphere is dominated by the water group ions: O+, H2O+, OH+ and others, hydronium ion (H3O+), HO2+ and O2+, although protons and nitrogen ions (N+) are also present. The main source of water is Enceladus, which releases 300–600 kg/s of water vapor from the geysers near its south pole. The released water and hydroxyl (OH) radicals (a product of water's dissociation) form a rather thick torus around the moon's orbit at 4 Rs with densities up to 10,000 molecules per cubic centimeter. At least 100 kg/s of this water is eventually ionized and added to the co–rotating magnetospheric plasma. Additional sources of water group ions are Saturn's rings and other icy moons. The Cassini spacecraft also observed small amounts of N+ ions in the inner magnetosphere, which probably originate from Enceladus as well. In the outer parts of the magnetosphere the dominant ions are protons, which originate either from the Solar wind or Saturn's ionosphere. Titan, which orbits close to the magnetopause boundary at 20 Rs, is not a significant source of plasma. The relatively cold plasma in the innermost region of Saturn's magnetosphere, inside 3 Rs (near the rings) consists mainly of O+ and O2+ ions. There ions together with electrons form an ionosphere surrounding the saturnian rings. For both Jupiter and Saturn, transport of plasma from the inner to the outer parts of the magnetosphere is thought to be related to interchange instability. In the case of Saturn, charge exchange facilitates the transfer of energy from the previously hot ions to the neutral gases in the inner magnetosphere. Then, magnetic flux tubes loaded with this newly cold, water–rich plasma interchange with flux tubes filled with hot plasma arriving from the outer magnetosphere. The instability is driven by centrifugal force exerted by the plasma on the magnetic field. The cold plasma is eventually removed from the magnetosphere by plasmoids formed when the magnetic field reconnects in the magnetotail. The plasmoids move down the tail and escape from the magnetosphere. The reconnection or substorm process is thought to be under the control of the solar wind and Saturn's largest moon Titan, which orbits near the outer boundary of the magnetosphere. In the magnetodisk region, beyond 6 Rs, the plasma within the co–rotating sheet exerts a significant centrifugal force on the magnetic field, causing it to stretch. This interaction creates a current in the equatorial plane flowing azimuthally with rotation and extending as far as 20 Rs from the planet. The total strength of this current varies from 8 to 17 MA. The ring current in the saturnian magnetosphere is highly variable and depends on the solar wind pressure, being stronger when the pressure is weaker. The magnetic moment associated with this current slightly (by about 10 nT) depresses the magnetic field in the inner magnetosphere, although it increases the total magnetic moment of the planet and causing the size of the magnetosphere to become larger. Aurorae Saturn has bright polar aurorae, which have been observed in the ultraviolet, visible and near infrared light. The aurorae usually look like bright continuous circles (ovals) surrounding the poles of the planet. The latitude of auroral ovals varies in the range of 70–80°; the average position is for the southern aurora, while the northern aurora is closer to the pole by about 1.5°. From time to time either aurorae can assume a spiral shape instead of oval. In this case it begins near midnight at a latitude of around 80°, then its latitude decreases to as low as 70° as it continues into the dawn and day sectors (counterclockwise). In the dusk sector the auroral latitude increases again, although when it returns to the night sector it still has a relatively low latitude and does not connect to the brighter dawn part. Unlike Jupiter's, Saturn's main auroral ovals are not related to the breakdown of the co–rotation of the plasma in the outer parts of the planet's magnetosphere. The aurorae on Saturn are thought to be connected to the reconnection of the magnetic field under the influence of the Solar wind (Dungey cycle), which drives an upward current (about 10 million amperes) from the ionosphere and leads to the acceleration and precipitation of energetic (1–10 keV) electrons into the polar thermosphere of Saturn. The Saturnian aurorae are more similar to those of the Earth, where they are also Solar wind driven. The ovals themselves correspond to the boundaries between open and closed magnetic field lines—so called polar caps, which are thought to reside at the distance of 10–15° from the poles. The aurorae of Saturn are highly variable. Their location and brightness strongly depends on the Solar wind pressure: the aurorae become brighter and move closer to the poles when the Solar wind pressure increases. The bright auroral features are observed to rotate with the angular speed of 60–75% that of Saturn. From time to time bright features appear in the dawn sector of the main oval or inside it. The average total power emitted by the aurorae is about 50 GW in the far ultraviolet (80–170 nm) and 150–300 GW in the near-infrared (3–4 μm—H3+ emissions) parts of the spectrum. Saturn kilometric radiation Saturn is the source of rather strong low frequency radio emissions called Saturn kilometric radiation (SKR). The frequency of SKR lies in the range 10–1300 kHz (wavelength of a few kilometers) with the maximum around 400 kHz. The power of these emissions is strongly modulated by the rotation of the planet and is correlated with changes in the solar wind pressure. For instance, when Saturn was immersed into the giant magnetotail of Jupiter during Voyager 2 flyby in 1981, the SKR power decreased greatly or even ceased completely. The kilometric radiation is thought to be generated by the Cyclotron Maser Instability of the electrons moving along magnetic field lines related to the auroral regions of Saturn. Thus the SKR is related to the auroras around the poles of the planet. The radiation itself comprises spectrally diffuse emissions as well as narrowband tones with bandwidths as narrow as 200 Hz. In the frequency–time plane, arc-like features are often observed, much like in the case of the Jovian kilometric radiation. The total power of the SKR is around 1 GW. The modulation of the radio emissions by planetary rotation is traditionally used to determine the rotation period of the interiors of fluid giant planets. In the case of Saturn, however, this appears to be impossible, as the period varies at the timescale of ten years. In 1980–1981 the periodicity in the radio emissions as measured by Voyager 1 and 2 was , which was then adopted as the rotational period of Saturn. Scientists were surprised when Galileo and then Cassini returned a different value—. Further observation indicated that the modulation period changes by as much as 1% on the characteristic timescale of 20–30 days with an additional long-term trend. There is a correlation between the period and solar wind speed, however, the causes of this change remain a mystery. One reason may be that the Saturnian perfectly axially symmetric magnetic field fails to impose a strict corotation on the magnetospheric plasma making it slip relative to the planet. The lack of a precise correlation between the variation period of SKR and planetary rotation makes it all but impossible to determine the true rotational period of Saturn. Radiation belts Saturn has relatively weak radiation belts, because energetic particles are absorbed by the moons and particulate material orbiting the planet. The densest (main) radiation belt lies between the inner edge of the Enceladus gas torus at 3.5 Rs and the outer edge of the A Ring at 2.3 Rs. It contains protons and relativistic electrons with energies from hundreds of kiloelectronvolts (keV) to as high as tens of megaelectronvolts (MeV) and possibly other ions. Beyond 3.5 Rs the energetic particles are absorbed by the neutral gas and their numbers drop, although less energetic particles with energies in the range of hundreds keV appear again beyond 6 Rs—these are the same particles that contribute to the ring current. The electrons in the main belt probably originate in the outer magnetosphere or Solar wind, from which they are transported by the diffusion and then adiabatically heated. However, the energetic protons consist of two populations of particles. The first population with energies of less than about 10 MeV has the same origin as electrons, while the second one with the maximum flux near 20 MeV results from the interaction of cosmic rays with solid material present in the Saturnian system (so called cosmic ray albedo neutron decay process—CRAND). The main radiation belt of Saturn is strongly influenced by interplanetary solar wind disturbances. The innermost region of the magnetosphere near the rings is generally devoid of energetic ions and electrons because they are absorbed by ring particles. Saturn, however, has the second radiation belt discovered by Cassini in 2004 and located just inside the innermost D Ring. This belt probably consists of energetic charged particles formed via the CRAND process or of ionized energetic neutral atoms coming from the main radiation belt. The saturnian radiation belts are generally much weaker than those of Jupiter and do not emit much microwave radiation (with frequency of a few Gigahertz). Estimates shows that their decimetric radio emissions (DIM) would be impossible to detect from the Earth. Nevertherless the high energy particles cause weathering of the surfaces of the icy moons and sputter water, water products and oxygen from them. Interaction with rings and moons The abundant population of solid bodies orbiting Saturn including moons as well as ring particles exerts a strong influence on the magnetosphere of Saturn. The plasma in the magnetosphere co-rotates with the planet, continuously impinging on the trailing hemispheres of slowly moving moons. While ring particles and the majority of moons only passively absorb plasma and energetic charged particles, three moons – Enceladus, Dione and Titan – are significant sources of new plasma. The absorption of energetic electrons and ions reveals itself by noticeable gaps in the radiation belts of Saturn near the moon's orbits, while the dense rings of Saturn eliminate all energetic electrons and ions closer than 2.2 RS, creating a low radiation zone in the vicinity of the planet. The absorption of the co-rotating plasma by a moon disturbs the magnetic field in its empty wake—the field is pulled towards a moon, creating a region of a stronger magnetic field in the near wake. The three moons mentioned above add new plasma into the magnetosphere. By far the strongest source is Enceladus, which ejects a fountain of water vapor, carbon dioxide and nitrogen through cracks in its south pole region. A fraction of this gas is ionized by the hot electrons and solar ultraviolet radiation and is added to the co-rotational plasma flow. Titan once was thought to be the principal source of plasma in Saturn's magnetosphere, especially of nitrogen. The new data obtained by Cassini in 2004–2008 established that it is not a significant source of nitrogen after all, although it may still provide significant amounts of hydrogen (due to dissociation of methane). Dione is the third moon producing more new plasma than it absorbs. The mass of plasma created in the vicinity of it (about 6 g/s) is about 1/300 as much as near Enceladus. However, even this low value can not be explained only by sputtering of its icy surface by energetic particles, which may indicate that Dione is endogenously active like Enceladus. The moons that create new plasma slow the motion of the co-rotating plasma in their vicinity, which leads to the pile-up of the magnetic field lines in front of them and weakening of the field in their wakes—the field drapes around them. This is the opposite to what is observed for the plasma-absorbing moons. The plasma and energetic particles present in the magnetosphere of Saturn, when absorbed by ring particles and moons, cause radiolysis of the water ice. Its products include ozone, hydrogen peroxide and molecular oxygen. The first one has been detected in the surfaces of Rhea and Dione, while the second is thought to be responsible for the steep spectral slopes of moons' reflectivities in the ultraviolet region. The oxygen produced by radiolysis forms tenuous atmospheres around rings and icy moons. The ring atmosphere was detected by Cassini for the first time in 2004. A fraction of the oxygen gets ionized, creating a small population of O2+ ions in the magnetosphere. The influence of Saturn's magnetosphere on its moons is more subtle than the influence of Jupiter on its moons. In the latter case, the magnetosphere contains a significant number of sulfur ions, which, when implanted in surfaces, produce characteristic spectral signatures. In the case of Saturn, the radiation levels are much lower and the plasma is composed mainly of water products, which, when implanted, are indistinguishable from the ice already present. Exploration As of 2014 the magnetosphere of Saturn has been directly explored by four spacecraft. The first mission to study the magnetosphere was Pioneer 11 in September 1979. Pioneer 11 discovered the magnetic field and made some measurements of the plasma parameters. In November 1980 and August 1981, Voyager 1–2 probes investigated the magnetosphere using an improved set of instruments. From the fly-by trajectories they measured the planetary magnetic field, plasma composition and density, high energy particle energy and spatial distribution, plasma waves and radio emissions. Cassini spacecraft was launched in 1997, and arrived in 2004, making the first measurements in more than two decades. The spacecraft continued to provide information about the magnetic field and plasma parameters of the saturnian magnetosphere until its intentional destruction on September 15, 2017. In the 1990s, the Ulysses spacecraft conducted extensive measurements of the Saturnian kilometric radiation (SKR), which is unobservable from Earth due to the absorption in the ionosphere. The SKR is powerful enough to be detected from a spacecraft at the distance of several astronomical units from the planet. Ulysses discovered that the period of the SKR varies by as much as 1%, and therefore is not directly related to the rotation period of the interior of Saturn.	Physical sciences	Solar System	Astronomy
2309752	https://en.wikipedia.org/wiki/Komatiite	Komatiite	Komatiite is a type of ultramafic mantle-derived volcanic rock defined as having crystallised from a lava of at least 18 wt% magnesium oxide (MgO). It is classified as a 'picritic rock'. Komatiites have low silicon, potassium and aluminium, and high to extremely high magnesium content. Komatiite was named for its type locality along the Komati River in South Africa, and frequently displays spinifex texture composed of large dendritic plates of olivine and pyroxene. Komatiites are rare rocks; almost all komatiites were formed during the Archaean Eon (4.03–2.5 billion years ago), with few younger (Proterozoic or Phanerozoic) examples known. This restriction in age is thought to be due to cooling of the mantle, which may have been hotter during the Archaean. The early Earth had much higher heat production, due to the residual heat from planetary accretion, as well as the greater abundance of radioactive isotopes, particularly shorter lived ones like uranium 235 which produce more decay heat. Lower temperature mantle melts such as basalt and picrite have essentially replaced komatiites as an eruptive lava on the Earth's surface. Geographically, komatiites are predominantly restricted in distribution to the Archaean shield areas, and occur with other ultramafic and high-magnesian mafic volcanic rocks in Archaean greenstone belts. The youngest komatiites are from the island of Gorgona on the Caribbean oceanic plateau off the Pacific coast of Colombia, and a rare example of Proterozoic komatiite is found in the Winnipegosis komatiite belt in Manitoba, Canada. Petrology Magmas of komatiitic compositions have a very high melting point, with calculated eruption temperatures up to, and possibly in excess of 1600 °C. Basaltic lavas normally have eruption temperatures of about 1100 to 1250 °C. The higher melting temperatures required to produce komatiite have been attributed to the presumed higher geothermal gradients in the Archaean Earth. Komatiitic lava was extremely fluid when it erupted (possessing the viscosity close to that of water but with the density of rock). Compared to the basaltic lava of the Hawaiian plume basalts at ~1200 °C, which flows the way treacle or honey does, the komatiitic lava would have flowed swiftly across the surface, leaving extremely thin lava flows (down to 10 mm thick). The major komatiitic sequences preserved in Archaean rocks are thus considered to be lava tubes, ponds of lava etc., where the komatiitic lava accumulated. Komatiite chemistry is different from that of basaltic and other common mantle-produced magmas, because of differences in degrees of partial melting. Komatiites are considered to have been formed by high degrees of partial melting, usually greater than 50%, and hence have high MgO with low K2O and other incompatible elements. There are two geochemical classes of komatiite; aluminium undepleted komatiite (AUDK) (also known as Group I komatiites) and aluminium depleted komatiite (ADK) (also known as Group II komatiites), defined by their Al2O3/TiO2 ratios. These two classes of komatiite are often assumed to represent a real petrological source difference between the two types related to depth of melt generation. Al-depleted komatiites have been modeled by melting experiments as being produced by high degrees of partial melting at high pressure where garnet in the source is not melted, whereas Al-undepleted komatiites are produced by high degrees of partial melts at lesser depth. However, recent studies of fluid inclusions in chrome spinels from the cumulate zones of komatiite flows have shown that a single komatiite flow can be derived from the mixing of parental magmas with a range of Al2O3/TiO2 ratios, calling into question this interpretation of the formations of the different komatiite groups. Komatiites probably form in extremely hot mantle plumes or in Archaean subduction zones. Boninite magmatism is similar to komatiite magmatism but is produced by fluid-fluxed melting above a subduction zone. Boninites with 10–18% MgO tend to have higher large-ion lithophile elements (LILE: Ba, Rb, Sr) than komatiites. Mineralogy The pristine volcanic mineralogy of komatiites is composed of forsteritic olivine (Fo90 and upwards), calcic and often chromian pyroxene, anorthite (An85 and upwards) and chromite. A considerable population of komatiite examples show a cumulate texture and morphology. The usual cumulate mineralogy is highly magnesium rich forsterite olivine, though chromian pyroxene cumulates are also possible (though rarer). Volcanic rocks rich in magnesium may be produced by accumulation of olivine phenocrysts in basalt melts of normal chemistry: an example is picrite. Part of the evidence that komatiites are not magnesium-rich simply because of cumulate olivine is textural: some contain spinifex texture, a texture attributable to rapid crystallization of the olivine in a thermal gradient in the upper part of a lava flow. "Spinifex" texture is named after the common name for the Australian grass Triodia, which grows in clumps with similar shapes. Another line of evidence is that the MgO content of olivines formed in komatiites is toward the nearly pure MgO forsterite composition, which can only be achieved in bulk by crystallisation of olivine from a highly magnesian melt. The rarely preserved flow top breccia and pillow margin zones in some komatiite flows are essentially volcanic glass, quenched in contact with overlying water or air. Because they are rapidly cooled, they represent the liquid composition of the komatiites, and thus record an anhydrous MgO content of up to 32% MgO. Some of the highest magnesian komatiites with clear textural preservation are those of the Barberton belt in South Africa, where liquids with up to 34% MgO can be inferred using bulk rock and olivine compositions. The mineralogy of a komatiite varies systematically through the typical stratigraphic section of a komatiite flow and reflects magmatic processes which komatiites are susceptible to during their eruption and cooling. The typical mineralogical variation is from a flow base composed of olivine cumulate, to a spinifex textured zone composed of bladed olivine and ideally a pyroxene spinifex zone and olivine-rich chill zone on the upper eruptive rind of the flow unit. Primary (magmatic) mineral species also encountered in komatiites include olivine, the pyroxenes augite, pigeonite and bronzite, plagioclase, chromite, ilmenite and rarely pargasitic amphibole. Secondary (metamorphic) minerals include serpentine, chlorite, amphibole, sodic plagioclase, quartz, iron oxides and rarely phlogopite, baddeleyite, and pyrope or hydrogrossular garnet. Metamorphism All known komatiites have been metamorphosed, therefore should technically be termed 'metakomatiite' though the prefix meta is inevitably assumed. Many komatiites are highly altered and serpentinized or carbonated from metamorphism and metasomatism. This results in significant changes to the mineralogy and the texture. Hydration vs. carbonation The metamorphic mineralogy of ultramafic rocks, particularly komatiites, is only partially controlled by composition. The character of the connate fluids which are present during low temperature metamorphism whether prograde or retrograde control the metamorphic assemblage of a metakomatiite (hereafter the prefix meta- is assumed). The factor controlling the mineral assemblage is the partial pressure of carbon dioxide within the metamorphic fluid, called the XCO2. If XCO2 is above 0.5, the metamorphic reactions favor formation of talc, magnesite (magnesium carbonate), and tremolite amphibole. These are classed as talc-carbonation reactions. Below XCO2 of 0.5, metamorphic reactions in the presence of water favor production of serpentinite. There are thus two main classes of metamorphic komatiite; carbonated and hydrated. Carbonated komatiites and peridotites form a series of rocks dominated by the minerals chlorite, talc, magnesite or dolomite and tremolite. Hydrated metamorphic rock assemblages are dominated by the minerals chlorite, serpentine-antigorite and brucite. Traces of talc, tremolite and dolomite may be present, as it is very rare that no carbon dioxide is present in metamorphic fluids. At higher metamorphic grades, anthophyllite, enstatite, olivine and diopside dominate as the rock mass dehydrates. Mineralogic variations in komatiite flow facies Komatiite tends to fractionate from high-magnesium compositions in the flow bases where olivine cumulates dominate, to lower magnesium compositions higher up in the flow. Thus, the current metamorphic mineralogy of a komatiite will reflect the chemistry, which in turn represents an inference as to its volcanological facies and stratigraphic position. Typical metamorphic mineralogy is tremolite-chlorite, or talc-chlorite mineralogy in the upper spinifex zones. The more magnesian-rich olivine-rich flow base facies tend to be free from tremolite and chlorite mineralogy and are dominated by either serpentine-brucite +/- anthophyllite if hydrated, or talc-magnesite if carbonated. The upper flow facies tend to be dominated by talc, chlorite, tremolite, and other magnesian amphiboles (anthophyllite, cummingtonite, gedrite, etc.). For example, the typical flow facies (see below) may have the following mineralogy; Geochemistry Komatiite can be classified according to the following IUGS geochemical criteria: SiO2 less than 52 wt% MgO more than 18 wt% K2O + Na2O less than 1 wt% TiO2 less than 1 wt% When meeting the above, but the TiO2 is more than 1 wt%, it is classified as meimechite. A similar high-Mg volcanic rock is boninite, having 52–63 wt% SiO2, more than 8 wt% MgO and less than 0.5 wt% TiO2. The above geochemical classification must be the essentially unaltered magma chemistry and not the result of crystal accumulation (as in peridotite). Through a typical komatiite flow sequence the chemistry of the rock will change according to the internal fractionation which occurs during eruption. This tends to lower MgO, Cr, Ni, and increase Al, K2O, Na, CaO and SiO2 toward the top of the flow. Rocks rich in MgO, K2O, Ba, Cs, and Rb may be lamprophyres, kimberlites or other rare ultramafic, potassic or ultrapotassic rocks. Morphology and occurrence Komatiites often show pillow lava structure, autobrecciated upper margins consistent with underwater eruption forming a rigid upper skin to the lava flows. Proximal volcanic facies are thinner and interleaved with sulfidic sediments, black shales, cherts and tholeiitic basalts. Komatiites were produced from a relatively wet mantle. Evidence of this is from their association with felsics, occurrences of komatiitic tuffs, niobium anomalies and by S- and H2O-borne rich mineralizations. Textural features A common and distinctive texture is known as spinifex texture and consists of long acicular phenocrysts of olivine (or pseudomorphs of alteration minerals after olivine) or pyroxene which give the rock a bladed appearance especially on a weathered surface. Spinifex texture is the result of rapid crystallization of highly magnesian liquid in the thermal gradient at the margin of the flow or sill. Harrisite texture, first described from intrusive rocks (not komatiites) at Harris Bay on the island of Rùm in Scotland, is formed by nucleation of crystals on the floor of a magma chamber. Harrisites are known to form megacrystal aggregates of pyroxene and olivine up to 1 metre in length. Harrisite texture is found in some very thick lava flows of komatiite, for example in the Norseman-Wiluna Greenstone Belt of Western Australia, in which crystallization of cumulates has occurred. Volcanology Komatiite volcano morphology is interpreted to have the general form and structure of a shield volcano, typical of most large basalt edifices, as the magmatic event which forms komatiites erupts less magnesian materials. However, the initial flux of the most magnesian magmas is interpreted to form a channelised flow facie, which is envisioned as a fissure vent releasing highly fluid komatiitic lava onto the surface. This then flows outwards from the vent fissure, concentrating into topographical lows, and forming channel environments composed of high MgO olivine adcumulate flanked by a 'sheeted flow facies' aprons of lower MgO olivine and pyroxene thin-flow spinifex sheets. The typical komatiite lava flow has six stratigraphically related elements; A1 – pillowed and variolitic chilled flow top, often grading and transitional with sediment A2 – Zone of quickly chilled, feathery acicular olivine-clinopyroxene-glass representing a chilled margin on the top of the flow unit A3 – Olivine spinifex sequence composed of sheaf and book-like olivine spinifex, representing a downward-growing crystal accumulation on the flow top B1 – Olivine mesocumulate to orthocumulate, representing a harrisite grown in flowing liquid melt B2 – Olivine adcumulate composed of >93% interlocking equant olivine crystals B3 – Lower chill margin composed of olivine adcumulate to mesocumulate, with finer grain size. Individual flow units may not be entirely preserved, as subsequent flow units may thermally erode the A zone spinifex flows. In the distal thin flow facies, B zones are poorly developed to absent, as not enough through-flowing liquid existed to grow the adcumulate. The channel and sheeted flows are then covered by high-magnesian basalts and tholeiitic basalts as the volcanic event evolves to less magnesian compositions. The subsequent magmatism, being higher silica melts, tends to form a more typical shield volcano architecture. Intrusive komatiites Komatiite magma is extremely dense and unlikely to reach the surface, being more likely to pool lower within the crust. Modern (post-2004) interpretations of some of the larger olivine adcumulate bodies in the Yilgarn craton have revealed that the majority of komatiite olivine adcumulate occurrences are likely to be subvolcanic to intrusive in nature. This is recognised at the Mt Keith nickel deposit where wall-rock intrusive textures and xenoliths of felsic country rocks have been recognised within the low-strain contacts. The previous interpretations of these large komatiite bodies was that they were "super channels" or reactivated channels, which grew to over 500 m in stratigraphic thickness during prolonged volcanism. These intrusions are considered to be channelised sills, formed by injection of komatiitic magma into the stratigraphy, and inflation of the magma chamber. Economic nickel-mineralised olivine adcumulate bodies may represent a form of sill-like conduit, where magma pools in a staging chamber before erupting onto the surface. Economic importance The economic importance of komatiite was first widely recognised in the early 1960s with the discovery of massive nickel sulfide mineralisation at Kambalda, Western Australia. Komatiite-hosted nickel-copper sulfide mineralisation today accounts for about 14% of the world's nickel production, mostly from Australia, Canada and South Africa. Komatiites are associated with nickel and gold deposits in Australia, Canada, South Africa and most recently in the Guiana shield of South America.	Physical sciences	Igneous rocks	Earth science
9038260	https://en.wikipedia.org/wiki/Phosphorus%20cycle	Phosphorus cycle	The phosphorus cycle is the biogeochemical cycle that involves the movement of phosphorus through the lithosphere, hydrosphere, and biosphere. Unlike many other biogeochemical cycles, the atmosphere does not play a significant role in the movement of phosphorus, because phosphorus and phosphorus-based materials do not enter the gaseous phase readily, as the main source of gaseous phosphorus, phosphine, is only produced in isolated and specific conditions. Therefore, the phosphorus cycle is primarily examined studying the movement of orthophosphate (PO4)3-, the form of phosphorus that is most commonly seen in the environment, through terrestrial and aquatic ecosystems. Living organisms require phosphorus, a vital component of DNA, RNA, ATP, etc., for their proper functioning. Phosphorus also enters in the composition of phospholipids present in cell membranes. Plants assimilate phosphorus as phosphate and incorporate it into organic compounds. In animals, inorganic phosphorus in the form of apatite () is also a key component of bones, teeth (tooth enamel), etc. On the land, phosphorus gradually becomes less available to plants over thousands of years, since it is slowly lost in runoff. Low concentration of phosphorus in soils reduces plant growth and slows soil microbial growth, as shown in studies of soil microbial biomass. Soil microorganisms act as both sinks and sources of available phosphorus in the biogeochemical cycle. Short-term transformation of phosphorus is chemical, biological, or microbiological. In the long-term global cycle, however, the major transfer is driven by tectonic movement over geologic time and weathering of phosphate containing rock such as apatite. Furthermore, phosphorus tends to be a limiting nutrient in aquatic ecosystems. However, as phosphorus enters aquatic ecosystems, it has the possibility to lead to over-production in the form of eutrophication, which can happen in both freshwater and saltwater environments. Human activities have caused major changes to the global phosphorus cycle primarily through the mining and subsequent transformation of phosphorus minerals for use in fertilizer and industrial products. Some phosphorus is also lost as effluent through the mining and industrial processes as well. Phosphorus in the environment Ecological function Phosphorus is an essential nutrient for plants and animals. Phosphorus is a limiting nutrient for aquatic organisms. Phosphorus forms parts of important life-sustaining molecules that are very common in the biosphere. Phosphorus does enter the atmosphere in very small amounts when dust containing phosphorus is dissolved in rainwater and sea spray, but the element mainly remains on land and in rock and soil minerals. Phosphates which are found in fertilizers, sewage and detergents, can cause pollution in lakes and streams. Over-enrichment of phosphate in both fresh and inshore marine waters can lead to massive algae blooms. In fresh water, the death and decay of these blooms leads to eutrophication. An example of this is the Canadian Experimental Lakes Area. Freshwater algal blooms are generally caused by excess phosphorus, while those that take place in saltwater tend to occur when excess nitrogen is added. However, it is possible for eutrophication to be due to a spike in phosphorus content in both freshwater and saltwater environments. Phosphorus occurs most abundantly in nature as part of the orthophosphate ion (PO4)3−, consisting of a P atom and 4 oxygen atoms. On land most phosphorus is found in rocks and minerals. Phosphorus-rich deposits have generally formed in the ocean or from guano, and over time, geologic processes bring ocean sediments to land. Weathering of rocks and minerals release phosphorus in a soluble form where it is taken up by plants, and it is transformed into organic compounds. The plants may then be consumed by herbivores and the phosphorus is either incorporated into their tissues or excreted. After death, the animal or plant decays, and phosphorus is returned to the soil where a large part of the phosphorus is transformed into insoluble compounds. Runoff may carry a small part of the phosphorus back to the ocean. Generally with time (thousands of years) soils become deficient in phosphorus leading to ecosystem retrogression. Major pools in aquatic systems There are four major pools of phosphorus in freshwater ecosystems: dissolved inorganic phosphorus (DIP), dissolved organic phosphorus (DOP), particulate inorganic phosphorus (PIP) and particulate organic phosphorus (POP). Dissolved material is defined as substances that pass through a 0.45 μm filter. DIP consists mainly of orthophosphate (PO43-) and polyphosphate, while DOP consists of DNA and phosphoproteins. Particulate matter are the substances that get caught on a 0.45 μm filter and do not pass through. POP consists of both living and dead organisms, while PIP mainly consists of hydroxyapatite, Ca5(PO4)3OH . Inorganic phosphorus comes in the form of readily soluble orthophosphate. Particulate organic phosphorus occurs in suspension in living and dead protoplasm and is insoluble. Dissolved organic phosphorus is derived from the particulate organic phosphorus by excretion and decomposition and is soluble. Biological function The primary biological importance of phosphates is as a component of nucleotides, which serve as energy storage within cells (ATP) or when linked together, form the nucleic acids DNA and RNA. The double helix of our DNA is only possible because of the phosphate ester bridge that binds the helix. Besides making biomolecules, phosphorus is also found in bone and the enamel of mammalian teeth, whose strength is derived from calcium phosphate in the form of hydroxyapatite. It is also found in the exoskeleton of insects, and phospholipids (found in all biological membranes). It also functions as a buffering agent in maintaining acid base homeostasis in the human body. Phosphorus cycling Phosphates move quickly through plants and animals; however, the processes that move them through the soil or ocean are very slow, making the phosphorus cycle overall one of the slowest biogeochemical cycles. The global phosphorus cycle includes four major processes: (i) tectonic uplift and exposure of phosphorus-bearing rocks such as apatite to surface weathering; (ii) physical erosion, and chemical and biological weathering of phosphorus-bearing rocks to provide dissolved and particulate phosphorus to soils, lakes and rivers; (iii) riverine and subsurface transportation of phosphorus to various lakes and run-off to the ocean; (iv) sedimentation of particulate phosphorus (e.g., phosphorus associated with organic matter and oxide/carbonate minerals) and eventually burial in marine sediments (this process can also occur in lakes and rivers). In terrestrial systems, bioavailable P (‘reactive P’) mainly comes from weathering of phosphorus-containing rocks. The most abundant primary phosphorus-mineral in the crust is apatite, which can be dissolved by natural acids generated by soil microbes and fungi, or by other chemical weathering reactions and physical erosion. The dissolved phosphorus is bioavailable to terrestrial organisms and plants and is returned to the soil after their decay. Phosphorus retention by soil minerals (e.g., adsorption onto iron and aluminum oxyhydroxides in acidic soils and precipitation onto calcite in neutral-to-calcareous soils) is usually viewed as the most important process in controlling terrestrial P-bioavailability in the mineral soil. This process can lead to the low level of dissolved phosphorus concentrations in soil solution. Various physiological strategies are used by plants and microorganisms for obtaining phosphorus from this low level of phosphorus concentration. Soil phosphorus is usually transported to rivers and lakes and can then either be buried in lake sediments or transported to the ocean via river runoff. Atmospheric phosphorus deposition is another important marine phosphorus source to the ocean. In surface seawater, dissolved inorganic phosphorus, mainly orthophosphate (PO43-), is assimilated by phytoplankton and transformed into organic phosphorus compounds. Phytoplankton cell lysis releases cellular dissolved inorganic and organic phosphorus to the surrounding environment. Some of the organic phosphorus compounds can be hydrolyzed by enzymes synthesized by bacteria and phytoplankton and subsequently assimilated. The vast majority of phosphorus is remineralized within the water column, and approximately 1% of associated phosphorus carried to the deep sea by the falling particles is removed from the ocean reservoir by burial in sediments. A series of diagenetic processes act to enrich sediment pore water phosphorus concentrations, resulting in an appreciable benthic return flux of phosphorus to overlying bottom waters. These processes include (i) microbial respiration of organic matter in sediments, (ii) microbial reduction and dissolution of iron and manganese (oxyhydr)oxides with subsequent release of associated phosphorus, which connects the phosphorus cycle to the iron cycle, and (iii) abiotic reduction of iron (oxyhydr)oxides by hydrogen sulfide and liberation of iron-associated phosphorus. Additionally, (iv) phosphate associated with calcium carbonate and (v) transformation of iron oxide-bound phosphorus to vivianite play critical roles in phosphorus burial in marine sediments. These processes are similar to phosphorus cycling in lakes and rivers. Although orthophosphate (PO43-), the dominant inorganic P species in nature, is oxidation state (P5+), certain microorganisms can use phosphonate and phosphite (P3+ oxidation state) as a P source by oxidizing it to orthophosphate. Recently, rapid production and release of reduced phosphorus compounds has provided new clues about the role of reduced P as a missing link in oceanic phosphorus. Phosphatic minerals The availability of phosphorus in an ecosystem is restricted by its rate of release during weathering. The release of phosphorus from apatite dissolution is a key control on ecosystem productivity. The primary mineral with significant phosphorus content, apatite [Ca5(PO4)3OH] undergoes carbonation. Little of this released phosphorus is taken up by biota, as it mainly reacts with other soil minerals. This leads to phosphorus becoming unavailable to organisms in the later stage of weathering and soil development as it will precipitate into rocks. Available phosphorus is found in a biogeochemical cycle in the upper soil profile, while phosphorus found at lower depths is primarily involved in geochemical reactions with secondary minerals. Plant growth depends on the rapid root uptake of phosphorus released from dead organic matter in the biochemical cycle. Phosphorus is limited in supply for plant growth. Phosphates move quickly through plants and animals; however, the processes that move them through the soil or ocean are very slow, making the phosphorus cycle overall one of the slowest biogeochemical cycles. Low-molecular-weight (LMW) organic acids are found in soils. They originate from the activities of various microorganisms in soils or may be exuded from the roots of living plants. Several of those organic acids are capable of forming stable organo-metal complexes with various metal ions found in soil solutions. As a result, these processes may lead to the release of inorganic phosphorus associated with aluminum, iron, and calcium in soil minerals. The production and release of oxalic acid by mycorrhizal fungi explain their importance in maintaining and supplying phosphorus to plants. The availability of organic phosphorus to support microbial, plant and animal growth depends on the rate of their degradation to generate free phosphate. There are various enzymes such as phosphatases, nucleases and phytase involved for the degradation. Some of the abiotic pathways in the environment studied are hydrolytic reactions and photolytic reactions. Enzymatic hydrolysis of organic phosphorus is an essential step in the biogeochemical phosphorus cycle, including the phosphorus nutrition of plants and microorganisms and the transfer of organic phosphorus from soil to bodies of water. Many organisms rely on the soil derived phosphorus for their phosphorus nutrition. Eutrophication Eutrophication is when waters are enriched by nutrients that lead to structural changes to the aquatic ecosystem such as algae bloom, deoxygenation, reduction of fish species. It does occur naturally, as when lakes age they become more productive due to increases in major limiting reagents such as nitrogen and phosphorus. For example, phosphorus can enter into lakes where it will accumulate in the sediments and the biosphere. It can also be recycled from the sediments and the water system allowing it to stay in the environment. Antrhopogenic effects can also cause phosphorus to flow into aquatic ecosystems as seen in drainage water and runoff from fertilized soils on agricultural land. Additionally, eroded soils, which can be caused by deforestation and urbanization, can lead to more phosphorus and nitrogen being added to these aquatic ecosystems. These all increase the amount of phosphorus that enters the cycle which has led to excessive nutrient intake in freshwater systems causing dramatic growth in algal populations. When these algae die, their putrefaction depletes the water of oxygen and can toxify the waters. Both these effects cause plant and animal death rates to increase as the plants take in and animals drink the poisonous water. Saltwater phosphorus eutrophication Oceanic ecosystems gather phosphorus through many sources, but it is mainly derived from weathering of rocks containing phosphorus which are then transported to the oceans in a dissolved form by river runoff. Due to a dramatic rise in mining for phosphorus, it is estimated that humans have increased the net storage of phosphorus in soil and ocean systems by 75%. This increase in phosphorus has led to more eutrophication in ocean waters as phytoplankton blooms have caused a drastic shift in anoxic conditions seen in both the Gulf of Mexico and the Baltic Sea. Some research suggests that when anoxic conditions arise from eutrophication due to excess phosphorus, this creates a positive feedback loop that releases more phosphorus from oceanic reserves, exacerbating the issue. This could possibly create a self-sustaining cycle of oceanic anoxia where the constant recovery of phosphorus keeps stabilizing the eutrophic growth. Attempts to mitigate this problem using biological approaches are being investigated. One such approach involves using phosphorus accumulating organisms such as, Candidatus accumulibacter phosphatis, which are capable of effectively storing phosphorus in the form of phosphate in marine ecosystems. Essentially, this would alter how the phosphorus cycle exists currently in marine ecosystems. Currently, there has been a major influx of phosphorus due to increased agricultural use and other industrial applications, thus these organisms could theoretically store phosphorus and hold on to it until it could be recycled in terrestrial ecosystems which would have lost this excess phosphorus due to runoff. Wetland Wetlands are frequently applied to solve the issue of eutrophication. Nitrate is transformed in wetlands to free nitrogen and discharged to the air. Phosphorus is adsorbed by wetland soils which are taken up by the plants. Therefore, wetlands could help to reduce the concentration of nitrogen and phosphorus to remit eutrophication. However, wetland soils can only hold a limited amount of phosphorus. To remove phosphorus continually, it is necessary to add more new soils within the wetland from remnant plant stems, leaves, root debris, and undecomposable parts of dead algae, bacteria, fungi, and invertebrates. Human influences Nutrients are important to the growth and survival of living organisms and, hence, are essential for developing and maintaining healthy ecosystems. Humans have greatly influenced the phosphorus cycle by mining phosphate rock. For millennia, phosphorus was primarily brought into the environment by weathering phosphate-containing rocks, which would replenish the phosphorus normally lost to the environment through processes such as runoff, albeit on a very slow and gradual time scale. Since the 1840s, when the technology to mine and extract phosphorus became more prevalent, approximately 110 teragrams of phosphorus has been added to the environment. This trend appears to be continuing in the future as from 1900-2022, the amount of phosphorus mined globally has increased 72-fold, with an expected annual increase of 4%. Most of this mining is done to produce fertilizers which can be used on a global scale. However, at the rate humans are mining, the geological system can not quickly restore what is lost. Thus, researchers are examining ways to better recycle phosphorus in the environment, with one promising application including the use of microorganisms. Regardless, humans have had a profound impact on the phosphorus cycle with wide-reaching implications about food security, eutrophication, and the overall availability of the nutrient. Other human processes can have detrimental effects on the phosphorus cycle, such as the repeated application of liquid hog manure in excess to crops. Applying biosolids may also increase available phosphorus in soil. In poorly drained soils or in areas where snowmelt can cause periodic waterlogging, reducing conditions can be attained in 7–10 days. This causes a sharp increase in phosphorus concentration in solution, and phosphorus can be leached. In addition, reducing the soil causes a shift in phosphorus from resilient to more labile forms. This could eventually increase the potential for phosphorus loss. This is of particular concern for the environmentally sound management of such areas, where disposal of agricultural wastes has already become a problem. It is suggested that soil water regimes used for organic waste disposal be considered when preparing waste management regulations.	Physical sciences	Earth science basics: General	Earth science
9040547	https://en.wikipedia.org/wiki/Human%20skin	Human skin	The human skin is the outer covering of the body and is the largest organ of the integumentary system. The skin has up to seven layers of ectodermal tissue guarding muscles, bones, ligaments and internal organs. Human skin is similar to most of the other mammals' skin, and it is very similar to pig skin. Though nearly all human skin is covered with hair follicles, it can appear hairless. There are two general types of skin: hairy and glabrous skin (hairless). The adjective cutaneous literally means "of the skin" (from Latin cutis, skin). Skin plays an important immunity role in protecting the body against pathogens and excessive water loss. Its other functions are insulation, temperature regulation, sensation, synthesis of vitamin D, and the protection of vitamin B folates. Severely damaged skin will try to heal by forming scar tissue. This is often discoloured and depigmented. In humans, skin pigmentation (affected by melanin) varies among populations, and skin type can range from dry to non-dry and from oily to non-oily. Such skin variety provides a rich and diverse habitat for the approximately one thousand species of bacteria from nineteen phyla which have been found on human skin. Structure Human skin shares anatomical, physiological, biochemical and immunological properties with other mammalian lines. Pig skin especially shares similar epidermal and dermal thickness ratios to human skin: pig and human skin share similar hair follicle and blood vessel patterns; biochemically the dermal collagen and elastin content is similar in pig and human skin; and pig skin and human skin have similar physical responses to various growth factors. Skin has mesodermal cells which produce pigmentation, such as melanin provided by melanocytes, which absorb some of the potentially dangerous ultraviolet radiation (UV) in sunlight. It contains DNA repair enzymes that help reverse UV damage. People lacking the genes for these enzymes have high rates of skin cancer. One form predominantly produced by UV light, malignant melanoma, is particularly invasive, causing it to spread quickly, and can often be deadly. Human skin pigmentation varies substantially between populations; this has led to the classification of people(s) on the basis of skin colour. In terms of surface area, the skin is the second largest organ in the human body (the inside of the small intestine is 15 to 20 times larger). For the average adult human, the skin has a surface area of . The thickness of the skin varies considerably over all parts of the body, and between men and women, and young and old. An example is the skin on the forearm, which is on average in males and in females. of skin holds 650 sweat glands, 20 blood vessels, 60,000 melanocytes, and more than 1,000 nerve endings. The average human skin cell is about in diameter, but there are variants. A skin cell usually ranges from , depending on a variety of factors. Skin is composed of three primary layers: the epidermis, the dermis and the hypodermis. Epidermis The epidermis, "epi" coming from the Greek language meaning "over" or "upon", is the outermost layer of the skin. It forms the waterproof, protective wrap over the body's surface, which also serves as a barrier to infection and is made up of stratified squamous epithelium with an underlying basal lamina. The epidermis contains no blood vessels, and cells in the deepest layers are nourished almost exclusively by diffused oxygen from the surrounding air and to a far lesser degree by blood capillaries extending to the outer layers of the dermis. The main type of cells that make up the epidermis are Merkel cells, keratinocytes, with melanocytes and Langerhans cells also present. The epidermis can be further subdivided into the following strata (beginning with the outermost layer): corneum, lucidum (only in palms of hands and bottoms of feet), granulosum, spinosum, and basale. Cells are formed through mitosis at the basale layer. The daughter cells (see cell division) move up the strata changing shape and composition as they die due to isolation from their blood source. The cytoplasm is released and the protein keratin is inserted. They eventually reach the corneum and slough off (desquamation). This process is called "keratinization". This keratinized layer of skin is responsible for keeping water in the body and keeping other harmful chemicals and pathogens out, making skin a natural barrier to infection. The epidermis contains no blood vessels and is nourished by diffusion from the dermis. The main type of cells that make up the epidermis are keratinocytes, melanocytes, Langerhans cells, and Merkel cells. The epidermis helps the skin regulate body temperature. Layers The skin has up to seven layers of ectodermal tissue and guards the underlying muscles, bones, ligaments and internal organs. The epidermis is divided into several layers, where cells are formed through mitosis at the innermost layers. They move up the strata changing shape and composition as they differentiate and become filled with keratin. After reaching the top layer stratum corneum they are eventually 'sloughed off', or desquamated. This process is called keratinization and takes place within weeks. It was previously believed that the stratum corneum was "a simple, biologically inactive, outer epidermal layer comprising a fibrillar lattice of dead keratin". It is now understood that this is not true, and that the stratum corneum should be considered to be a live tissue. While it is true that the stratum corneum is mainly composed of terminally differentiated keratinocytes called corneocytes that are anucleated, these cells remain alive and metabolically functional until desquamated. Sublayers The epidermis is divided into the following 5 sublayers or strata: Stratum corneum Stratum lucidum Stratum granulosum Stratum spinosum Stratum basale (also called "stratum germinativum") Blood capillaries are found beneath the epidermis and are linked to an arteriole and a venule. Arterial shunt vessels may bypass the network in ears, the nose and fingertips. Genes and proteins expressed in the epidermis About 70% of all human protein-coding genes are expressed in the skin. Almost 500 genes have an elevated pattern of expression in the skin. There are fewer than 100 genes that are specific for the skin, and these are expressed in the epidermis. An analysis of the corresponding proteins show that these are mainly expressed in keratinocytes and have functions related to squamous differentiation and cornification. Dermis The dermis is the layer of skin beneath the epidermis that consists of connective tissue and cushions the body from stress and strain. The dermis is tightly connected to the epidermis by a basement membrane. It also harbours many nerve endings that provide the sense of touch and heat. It contains the hair follicles, sweat glands, sebaceous glands, apocrine glands, lymphatic vessels and blood vessels. The blood vessels in the dermis provide nourishment and waste removal from its own cells as well as from the stratum basale of the epidermis. The dermis is structurally divided into two areas: a superficial area adjacent to the epidermis, called the papillary region, and a deep thicker area known as the reticular region. Papillary region The papillary region is composed of loose areolar connective tissue. It is named for its finger-like projections called papillae, which extend toward the epidermis. The papillae provide the dermis with a "bumpy" surface that interdigitates with the epidermis, strengthening the connection between the two layers of skin. In the palms, fingers, soles, and toes, the influence of the papillae projecting into the epidermis forms contours in the skin's surface. These epidermal ridges occur in patterns (see: fingerprint) that are genetically and epigenetically determined and are therefore unique to the individual, making it possible to use fingerprints or footprints as a means of identification. Reticular region The reticular region lies deep in the papillary region and is usually much thicker. It is composed of dense irregular connective tissue, and receives its name from the dense concentration of collagenous, elastic, and reticular fibres that weave throughout it. These protein fibres give the dermis its properties of strength, extensibility, and elasticity. Also located within the reticular region are the roots of the hairs, sebaceous glands, sweat glands, receptors, nails, and blood vessels. Tattoo ink is held in the dermis. Stretch marks, often from pregnancy and obesity, are also located in the dermis. Subcutaneous tissue The subcutaneous tissue (also hypodermis and subcutis) is not part of the skin, but lies below the dermis of the cutis. Its purpose is to attach the skin to underlying bone and muscle as well as supplying it with blood vessels and nerves. It consists of loose connective tissue, adipose tissue and elastin. The main cell types are fibroblasts, macrophages and adipocytes (subcutaneous tissue contains 50% of body fat). Fat serves as padding and insulation for the body. Cross-section Cell count and cell mass Skin cell table The below table identifies the skin cell count and aggregate cell mass estimates for a 70 kg adult male (ICRP-23; ICRP-89, ICRP-110). Tissue mass is defined at 3.3 kg (ICRP-89, ICRP110) and addresses the skin's epidermis, dermis, hair follicles, and glands. The cell data is extracted from 'The Human Cell Count and Cell Size Distribution', Tissue-Table tab in the Supporting Information SO1 Dataset (xlsx). The 1200 record Dataset is supported by extensive references for cell size, cell count, and aggregate cell mass. Detailed data for below cell groups are further subdivided into all the cell types listed in the above sections and categorized by epidermal, dermal, hair follicle, and glandular subcategories in the dataset and on the dataset's graphical website interface. While adipocytes in the hypodermal adipose tissue are treated separately in the ICRP tissue categories, fat content (minus cell-membrane-lipids) resident in the dermal layer (Table-105, ICRP-23) is addressed by the below interstitial-adipocytes in the dermal layer. Development Skin colour Human skin shows high skin colour variety from the darkest brown to the lightest pinkish-white hues. Human skin shows higher variation in colour than any other single mammalian species and is the result of natural selection. Skin pigmentation in humans evolved to primarily regulate the amount of ultraviolet radiation (UVR) penetrating the skin, controlling its biochemical effects. The actual skin colour of different humans is affected by many substances, although the single most important substance determining human skin colour is the pigment melanin. Melanin is produced within the skin in cells called melanocytes and it is the main determinant of the skin colour of darker-skinned humans. The skin colour of people with light skin is determined mainly by the bluish-white connective tissue under the dermis and by the haemoglobin circulating in the veins of the dermis. The red colour underlying the skin becomes more visible, especially in the face, when, as consequence of physical exercise or the stimulation of the nervous system (anger, fear), arterioles dilate. There are at least five different pigments that determine the colour of the skin. These pigments are present at different levels and places. Melanin: It is brown in colour and present in the basal layer of the epidermis. Melanoid: It resembles melanin but is present diffusely throughout the epidermis. Carotene: This pigment is yellow to orange in colour. It is present in the stratum corneum and fat cells of dermis and superficial fascia. Hemoglobin (also spelled haemoglobin): It is found in blood and is not a pigment of the skin but develops a purple colour. Oxyhemoglobin: It is also found in blood and is not a pigment of the skin. It develops a red colour. There is a correlation between the geographic distribution of UV radiation (UVR) and the distribution of indigenous skin pigmentation around the world. Areas that highlight higher amounts of UVR reflect darker-skinned populations, generally located nearer towards the equator. Areas that are far from the tropics and closer to the poles have lower concentration of UVR, which is reflected in lighter-skinned populations. In the same population it has been observed that adult human females are considerably lighter in skin pigmentation than males. Females need more calcium during pregnancy and lactation, and vitamin D, which is synthesized from sunlight, helps in absorbing calcium. For this reason it is thought that females may have evolved to have lighter skin in order to help their bodies absorb more calcium. The Fitzpatrick scale is a numerical classification schema for human skin colour developed in 1975 as a way to classify the typical response of different types of skin to ultraviolet (UV) light: Ageing As skin ages, it becomes thinner and more easily damaged. Intensifying this effect is the decreasing ability of skin to heal itself as a person ages. Among other things, skin ageing is noted by a decrease in volume and elasticity. There are many internal and external causes to skin ageing. For example, ageing skin receives less blood flow and lower glandular activity. A validated comprehensive grading scale has categorized the clinical findings of skin ageing as laxity (sagging), rhytids (wrinkles), and the various facets of photoageing, including erythema (redness), and telangiectasia, dyspigmentation (brown discolouration), solar elastosis (yellowing), keratoses (abnormal growths) and poor texture. Cortisol causes degradation of collagen, accelerating skin ageing. Anti-ageing supplements are used to treat skin ageing. Photoageing Photoageing has two main concerns: an increased risk for skin cancer and the appearance of damaged skin. In younger skin, sun damage will heal faster since the cells in the epidermis have a faster turnover rate, while in the older population the skin becomes thinner and the epidermis turnover rate for cell repair is lower, which may result in the dermis layer being damaged. UV-induced DNA damage UV-irradiation of human skin cells generates damages in DNA through direct photochemical reactions at adjacent thymine or cytosine residues on the same strand of DNA. Cyclobutane pyrimidine dimers formed by two adjacent thymine bases, or by two adjacent cytosine bases, in DNA are the most frequent types of DNA damage induced by UV. Humans, as well as other organisms, are capable of repairing such UV-induced damages by the process of nucleotide excision repair. In humans this repair process protects against skin cancer. Types Though most human skin is covered with hair follicles, some parts can be hairless. There are two general types of skin, hairy and glabrous skin (hairless). The adjective cutaneous means "of the skin" (from Latin cutis, skin). Functions Skin performs the following functions: Protection: an anatomical barrier from pathogens and damage between the internal and external environment in bodily defence; Langerhans cells in the skin are part of the adaptive immune system. Perspiration contains lysozyme that break the bonds within the cell walls of bacteria. Sensation: contains a variety of nerve endings that react to heat and cold, touch, pressure, vibration, and tissue injury; see somatosensory system and haptics. Heat regulation: the skin contains a blood supply far greater than its requirements, which allows precise control of energy loss by radiation, convection and conduction. Dilated blood vessels increase perfusion and heat loss, while constricted vessels greatly reduce cutaneous blood flow and conserve heat. Control of evaporation: the skin provides a relatively dry and semi-impermeable barrier to fluid loss. Loss of this function contributes to the massive fluid loss in burns. Aesthetics and communication: others see our skin and can assess our mood, physical state and attractiveness. Storage and synthesis: acts as a storage centre for lipids and water, as well as a means of synthesis of vitamin D by action of UV on certain parts of the skin. Excretion: sweat contains urea, however its concentration is 1/130th that of urine, hence excretion by sweating is at most a secondary function to temperature regulation. Absorption: the cells comprising the outermost 0.25–0.40 mm of the skin are "almost exclusively supplied by external oxygen", although the "contribution to total respiration is negligible". In addition, medicine can be administered through the skin, by ointments or by means of adhesive patch, such as the nicotine patch or iontophoresis. The skin is an important site of transport in many other organisms. Water resistance: The skin acts as a water-resistant barrier so essential nutrients are not washed out of the body. Skin flora The human skin is a rich environment for microbes. Around 1,000 species of bacteria from 19 bacterial phyla have been found. Most come from only four phyla: Actinomycetota (51.8%), Bacillota (24.4%), Pseudomonadota (16.5%), and Bacteroidota (6.3%). Propionibacteria and Staphylococci species were the main species in sebaceous areas. There are three main ecological areas: moist, dry and sebaceous. In moist places on the body Corynebacteria together with Staphylococci dominate. In dry areas, there is a mixture of species but dominated by Betaproteobacteria and Flavobacteriales. Ecologically, sebaceous areas had greater species richness than moist and dry ones. The areas with least similarity between people in species were the spaces between fingers, the spaces between toes, axillae, and umbilical cord stump. Most similarly were beside the nostril, nares (inside the nostril), and on the back. Reflecting upon the diversity of the human skin researchers on the human skin microbiome have observed: "hairy, moist underarms lie a short distance from smooth dry forearms, but these two niches are likely as ecologically dissimilar as rainforests are to deserts." The NIH conducted the Human Microbiome Project to characterize the human microbiota, which includes that on the skin and the role of this microbiome in health and disease. Microorganisms like Staphylococcus epidermidis colonize the skin surface. The density of skin flora depends on region of the skin. The disinfected skin surface gets recolonized from bacteria residing in the deeper areas of the hair follicle, gut and urogenital openings. Clinical significance Diseases of the skin include skin infections and skin neoplasms (including skin cancer). Dermatology is the branch of medicine that deals with conditions of the skin. There are seven cervical, twelve thoracic, five lumbar, and five sacral. Certain diseases like shingles, caused by varicella-zoster infection, have pain sensations and eruptive rashes involving dermatomal distribution. Dermatomes are helpful in the diagnosis of vertebral spinal injury levels. Aside from the dermatomes, the epidermis cells are susceptible to neoplastic changes, resulting in various cancer types. The skin is also valuable for diagnosis of other conditions, since many medical signs show through the skin. Skin color affects the visibility of these signs, a source of misdiagnosis in unaware medical personnel. Society and culture Hygiene and skin care The skin supports its own ecosystems of microorganisms, including yeasts and bacteria, which cannot be removed by any amount of cleaning. Estimates place the number of individual bacteria on the surface of human skin at , though this figure varies greatly over the average of human skin. Oily surfaces, such as the face, may contain over . Despite these vast quantities, all of the bacteria found on the skin's surface would fit into a volume the size of a pea. In general, the microorganisms keep one another in check and are part of a healthy skin. When the balance is disturbed, there may be an overgrowth and infection, such as when antibiotics kill microbes, resulting in an overgrowth of yeast. The skin is continuous with the inner epithelial lining of the body at the orifices, each of which supports its own complement of microbes. Cosmetics should be used carefully on the skin because these may cause allergic reactions. Each season requires suitable clothing in order to facilitate the evaporation of the sweat. Sunlight, water and air play an important role in keeping the skin healthy. Oily skin Oily skin is caused by over-active sebaceous glands, that produce a substance called sebum, a naturally healthy skin lubricant. A high glycemic-index diet and dairy products (except for cheese) consumption increase IGF-1 generation, which in turn increases sebum production. Overwashing the skin does not cause sebum overproduction but may cause dryness. When the skin produces excessive sebum, it becomes heavy and thick in texture, known as oily skin. Oily skin is typified by shininess, blemishes and pimples. The oily-skin type is not necessarily bad, since such skin is less prone to wrinkling, or other signs of ageing, because the oil helps to keep needed moisture locked into the epidermis (outermost layer of skin). The negative aspect of the oily-skin type is that oily complexions are especially susceptible to clogged pores, blackheads, and buildup of dead skin cells on the surface of the skin. Oily skin can be sallow and rough in texture and tends to have large, clearly visible pores everywhere, except around the eyes and neck. Permeability Human skin has a low permeability; that is, most foreign substances are unable to penetrate and diffuse through the skin. Skin's outermost layer, the stratum corneum, is an effective barrier to most inorganic nanosized particles. This protects the body from external particles such as toxins by not allowing them to come into contact with internal tissues. However, in some cases it is desirable to allow particles entry to the body through the skin. Potential medical applications of such particle transfer has prompted developments in nanomedicine and biology to increase skin permeability. One application of transcutaneous particle delivery could be to locate and treat cancer. Nanomedical researchers seek to target the epidermis and other layers of active cell division where nanoparticles can interact directly with cells that have lost their growth-control mechanisms (cancer cells). Such direct interaction could be used to more accurately diagnose properties of specific tumours or to treat them by delivering drugs with cellular specificity. Nanoparticles Nanoparticles 40 nm in diameter and smaller have been successful in penetrating the skin. Research confirms that nanoparticles larger than 40 nm do not penetrate the skin past the stratum corneum. Most particles that do penetrate will diffuse through skin cells, but some will travel down hair follicles and reach the dermis layer. The permeability of skin relative to different shapes of nanoparticles has also been studied. Research has shown that spherical particles have a better ability to penetrate the skin compared to oblong (ellipsoidal) particles because spheres are symmetric in all three spatial dimensions. One study compared the two shapes and recorded data that showed spherical particles located deep in the epidermis and dermis whereas ellipsoidal particles were mainly found in the stratum corneum and epidermal layers. Nanorods are used in experiments because of their unique fluorescent properties but have shown mediocre penetration. Nanoparticles of different materials have shown skin's permeability limitations. In many experiments, gold nanoparticles 40 nm in diameter or smaller are used and have shown to penetrate to the epidermis. Titanium oxide (TiO2), zinc oxide (ZnO), and silver nanoparticles are ineffective in penetrating the skin past the stratum corneum. Cadmium selenide (CdSe) quantum dots have proven to penetrate very effectively when they have certain properties. Because CdSe is toxic to living organisms, the particle must be covered in a surface group. An experiment comparing the permeability of quantum dots coated in polyethylene glycol (PEG), PEG-amine, and carboxylic acid concluded the PEG and PEG-amine surface groups allowed for the greatest penetration of particles. The carboxylic acid coated particles did not penetrate past the stratum corneum. Increasing permeability Scientists previously believed that the skin was an effective barrier to inorganic particles. Damage from mechanical stressors was believed to be the only way to increase its permeability. Recently, simpler and more effective methods for increasing skin permeability have been developed. Ultraviolet radiation (UVR) slightly damages the surface of skin and causes a time-dependent defect allowing easier penetration of nanoparticles. The UVR's high energy causes a restructuring of cells, weakening the boundary between the stratum corneum and the epidermal layer. The damage of the skin is typically measured by the transepidermal water loss (TEWL), though it may take 3–5 days for the TEWL to reach its peak value. When the TEWL reaches its highest value, the maximum density of nanoparticles is able to permeate the skin. While the effect of increased permeability after UVR exposure can lead to an increase in the number of particles that permeate the skin, the specific permeability of skin after UVR exposure relative to particles of different sizes and materials has not been determined. There are other methods to increase nanoparticle penetration by skin damage: tape stripping is the process in which tape is applied to skin then lifted to remove the top layer of skin; skin abrasion is done by shaving the top 5–10 μm off the surface of the skin; chemical enhancement applies chemicals such as polyvinylpyrrolidone (PVP), dimethyl sulfoxide (DMSO), and oleic acid to the surface of the skin to increase permeability; electroporation increases skin permeability by the application of short pulses of electric fields. The pulses are high voltage and on the order of milliseconds when applied. Charged molecules penetrate the skin more frequently than neutral molecules after the skin has been exposed to electric field pulses. Results have shown molecules on the order of 100 μm to easily permeate electroporated skin. Applications A large area of interest in nanomedicine is the transdermal patch because of the possibility of a painless application of therapeutic agents with very few side effects. Transdermal patches have been limited to administer a small number of drugs, such as nicotine, because of the limitations in permeability of the skin. Development of techniques that increase skin permeability has led to more drugs that can be applied via transdermal patches and more options for patients. Increasing the permeability of skin allows nanoparticles to penetrate and target cancer cells. Nanoparticles along with multi-modal imaging techniques have been used as a way to diagnose cancer non-invasively. Skin with high permeability allowed quantum dots with an antibody attached to the surface for active targeting to successfully penetrate and identify cancerous tumours in mice. Tumour targeting is beneficial because the particles can be excited using fluorescence microscopy and emit light energy and heat that will destroy cancer cells. Sunblock and sunscreen Sunblock and sunscreen are different important skin-care products though both offer full protection from the sun. Sunblock—Sunblock is opaque and stronger than sunscreen, since it is able to block most of the UVA/UVB rays and radiation from the sun, and does not need to be reapplied several times in a day. Titanium dioxide and zinc oxide are two of the important ingredients in sunblock. Sunscreen—Sunscreen is more transparent once applied to the skin and also has the ability to protect against UVA/UVB rays, although the sunscreen's ingredients have the ability to break down at a faster rate once exposed to sunlight, and some of the radiation is able to penetrate to the skin. In order for sunscreen to be more effective it is necessary to consistently reapply and use one with a higher sun protection factor. Diet Vitamin A, also known as retinoids, benefits the skin by normalizing keratinization, downregulating sebum production, which contributes to acne, and reversing and treating photodamage, striae, and cellulite. Vitamin D and analogues are used to downregulate the cutaneous immune system and epithelial proliferation while promoting differentiation. Vitamin C is an antioxidant that regulates collagen synthesis, forms barrier lipids, regenerates vitamin E, and provides photoprotection. Vitamin E is a membrane antioxidant that protects against oxidative damage and also provides protection against harmful UV rays. Several scientific studies confirmed that changes in baseline nutritional status affects skin condition. Mayo Clinic lists foods they state help the skin: fruits and vegetables, whole-grains, dark leafy greens, nuts, and seeds.	Biology and health sciences	Human anatomy	Health
1078549	https://en.wikipedia.org/wiki/Aliivibrio%20fischeri	Aliivibrio fischeri	Aliivibrio fischeri (formerly Vibrio fischeri) is a Gram-negative, rod-shaped bacterium found globally in marine environments. This bacterium grows most effectively in water with a salt concentration at around 20g/L, and at temperatures between 24 and 28°C. This species is non-pathogenic and has bioluminescent properties. It is found predominantly in symbiosis with various marine animals, such as the Hawaiian bobtail squid. It is heterotrophic, oxidase-positive, and motile by means of a tuft of polar flagella. Free-living A. fischeri cells survive on decaying organic matter. The bacterium is a key research organism for examination of microbial bioluminescence, quorum sensing, and bacterial-animal symbiosis. It is named after Bernhard Fischer, a German microbiologist. Aliivibrio fischeri is the family Vibrionaceae. This family of bacteria tend to have adaptable metabolisms that can adjust to diverse circumstances. This flexibility may contribute to A. fischeri's ability to survive both alone and in symbiotic relationships. Ribosomal RNA comparison led to the reclassification of this species from genus Vibrio to the newly created Aliivibrio in 2007. The change is recognized as a valid publication, and according to the List of Prokaryotic names with Standing in Nomenclature (LPSN), the correct name. However, the name change is has not been universally adopted by most researchers, who still publish using the name Vibrio fischeri. Genome The genome of A. fischeri was completely sequenced in 2004 and consists of two chromosomes, one smaller and one larger. Chromosome 1 has 2.9 million base pairs (Mbp) and chromosome 2 has 1.5 Mbp, bringing the total genome to 4.4 Mbp. A. fischeri has the lowest G+C content of 27 Vibrio species but is still related to higher-pathogenicity species such as V. cholerae. The genome for A. fischeri also carries mobile genetic elements. The precise functions of these elements in A. fischeri are not fully understood. However, they are known to acquire new genes that are associated with virulence and resistance to environmental stresses in other bacterial genomes. Some strains of A. fischeri, such as strain ES114, contain a plasmid. The plasmid in strain ES114 is called pES100 and is most likely used for conjugation purposes. This purpose was determined based on the 45.8 kbp gene sequence, most of which codes for a type IV section system. The ability to preform conjugation can be helpful for both beneficial and pathogenic strains, as it allows for DNA exchange. There is evidence that the genome of A. fischeri includes pilus gene clusters. These clusters encode for many different kinds of pili, which serve a variety of functions. In this species, there are pili used for pathogenesis, twitching motility, tight adhesion, and toxin-coregulation, and more. Ecology A. fischeri are globally distributed in temperate and subtropical marine environments. They can be found free-floating in oceans, as well as associated with marine animals, sediment, and decaying matter. A. fischeri have been most studied as symbionts of marine animals, including squids in the genus Euprymna and Sepiola, where A. fischeri can be found in the squids' light organs. This relationship has been best characterized in the Hawaiian bobtail squid (Euprymna scolopes). A. fischeri is the only species of bacteria inhabiting the squid's light organ, despite an environment full of other bacteria. Symbiosis with the Hawaiian bobtail squid A. fischeri colonization of the light organ of the Hawaiian bobtail squid (Euprymna scolopes) is currently studied as a simple model for mutualistic symbiosis, as it contains only two species and A. fischeri can be cultured in a lab and genetically modified. Aliivibrio fischeri utilizes chitin as a primary carbon and nitrogen source in its symbiosis with the Hawaiian bobtail squid. In the squid’s light organ, A. fischeri breaks down chitin into N-acetylglucosamine (GlcNAc), which acts as both a nutrient and a chemoattractant, guiding colonization. Chitinases facilitate this breakdown, while the regulatory protein NagC controls gene expression for chitin and GlcNAc use. The bacteria metabolize GlcNAc through fermentation or respiration, supporting energy needs and bioluminescence, which are crucial for the mutualistic relationship with the squid. This mutualistic symbiosis provides A. fischeri with nutrients and a protected environment and helps the squid avoid predation using bioluminescence. A. fischeri provides luminescence by colonizing the light organ of the Hawaiian bobtail squid, which is on its ventral side. The organ luminesces at night, providing the squid with counter-illumination camouflage. The light organs of some squid contain reflective plates that intensify and direct the light produced, due to proteins known as reflectins. They regulate the light intensity to match that of the sea surface below. This strategy prevents the squid from casting a shadow on the ocean floor, helping it avoid predation during feeding. The A. fischeri population is maintained by daily cycles. About 90% of A. fischeri are ejected by the squid every morning in a process known as "venting". The 10% of bacteria remaining in the squid replenish the bacterial population before the following night. A. fischeri are horizontally acquired by young squids from their environment. Venting is thought to provide the source from which newly hatched squid are colonized. This colonization induces developmental and morphological changes in the squid's light organ, which is translucent. Morphological changes made by A. fischeri do not occur when the microbe cannot luminesce, such as a decrease in the number of pores in the light organ. Additionally, if colonization by A. fischeri is abruptly removed by antibiotics, the ciliated epithelium of the light organ will regress. These changes show that bioluminescence is truly essential for symbiosis. In the process of colonization, ciliated cells within the animals' photophores (light-producing organs) selectively draw in the symbiotic bacteria. These cells create microcurrents that, when combined with mucus, promote the growth of the symbionts and actively reject any competitors. The bacteria cause the ciliated cells to die once the light organ is sufficiently colonized. Bioluminescence The bioluminescence of A. fischeri is caused by transcription of the lux operon, and the following translation of the lux proteins, which produce the light. This process is induced through population-dependent quorum sensing. The population of A. fischeri needs to reach an optimal level to activate the lux operon and stimulate light production. The circadian rhythm controls light expression, where luminescence is much brighter during the day and dimmer at night, as required for camouflage. The bacterial luciferin-luciferase system is encoded by a set of genes labelled the lux operon. In A. fischeri, five such genes (luxCDABEG) have been identified as active in the emission of visible light, and two genes (luxR and luxI) are involved in regulating the operon. Several external and intrinsic factors appear to either induce or inhibit the transcription of this gene set and produce or suppress light emission. A. fischeri is one of many species of bacteria that commonly form symbiotic relationships with marine organisms. Marine organisms contain bacteria that use bioluminescence so they can find mates, ward off predators, attract prey, or communicate with other organisms. In return, the organism the bacteria are living within provides the bacteria with a nutrient-rich environment. The lux operon is a 9-kilobase fragment of the A. fischeri genome that controls bioluminescence through the catalytic activity of the enzyme luciferase. This operon has a known gene sequence of luxCDAB(F)E, where luxA and luxB code for the protein subunits of the luciferase enzyme, and the luxCDE codes for a fatty acid reductase complex that makes the fatty acids necessary for the luciferase mechanism. luxC codes for the enzyme acyl-reductase, luxD codes for acyl-transferase, and luxE makes the proteins needed for the enzyme acyl-protein synthetase. Luciferase produces blue/green light through the oxidation of reduced flavin mononucleotide and a long-chain aldehyde by diatomic oxygen. The reaction is summarized as: FMNH2 + O2 + R-CHO → FMN + R-COOH + H2O + light. The reduced flavin mononucleotide (FMNH) is provided by the fre gene, also referred to as luxG. In A. fischeri, it is directly next to luxE (giving luxCDABE-fre) from 1042306 to 1048745. To generate the aldehyde needed in the reaction above, three additional enzymes are needed. The fatty acids needed for the reaction are pulled from the fatty acid biosynthesis pathway by acyl-transferase. Acyl-transferase reacts with acyl-ACP to release R-COOH, a free fatty acid. R-COOH is reduced by a two-enzyme system to an aldehyde. The reaction is: R-COOH + ATP + NADPH → R-CHO + AMP + PP + NADP+. Quorum sensing One primary system that controls bioluminescence through regulation of the lux operon is quorum sensing, a conserved mechanism across many microbial species that regulates gene expression in response to bacterial concentration. Quorum sensing functions through the production of an autoinducer, usually a small organic molecule, by individual cells. As cell populations increase, levels of autoinducers increase, and specific proteins that regulate transcription of genes bind to these autoinducers, altering gene expression. This system allows microbial cells to "communicate" amongst each other and coordinate behaviors, such as luminescence, which require large amounts of cells to produce a noticeable effect. In A. fischeri, there are two primary quorum sensing systems, each of which responds to slightly different environments. The first system is commonly referred to as the lux system, as it is encoded within the lux operon, and uses the autoinducer 3OC6-HSL. The protein LuxI synthesizes this signal, which is subsequently released from the cell. This signal, 3OC6-HSL, then binds to the protein LuxR, which regulates the expression of many different genes, but most notably upregulation of genes involved in luminescence. The second system, commonly referred to as the ain system, uses the autoinducer C8-HSL, which is produced by the protein AinS. Similar to the lux system, the autoinducer C8-HSL increases activation of LuxR. In addition, C8-HSL binds to another transcriptional regulator, LitR, giving the ain and lux systems of quorum sensing slightly different genetic targets within the cell. The different genetic targets of the ain and lux systems are essential, because these two systems respond to different cellular environments. The ain system regulates transcription in response to intermediate cell density cell environments, producing lower levels of luminescence and even regulating metabolic processes such as the acetate switch. In contrast, the lux quorum sensing system occurs in response to high cell densities, producing high levels of luminescence and regulating the transcription of additional genes, including QsrP, RibB, and AcfA. Both of the ain and lux quorum sensing systems are essential for colonization of the squid and regulate multiple colonization factors in the bacteria. Research applications A. fischeri has broad applications in ecotoxicology and environmental research. Its bioluminescence is observed in oxygen-rich environments and thus is sensitive to toxicants. Reductions in light emissions are used in bioassays such as the Microtox test to assess water quality. It plays a key role in studying the effects of chemical mixtures, helping identify synergistic or antagonistic toxic interactions. In biotechnology, its light-producing mechanism is harnessed for developing biosensors that detect environmental pollutants in real time, making it a valuable tool in pollution monitoring and water treatment studies. Bioluminescence inhibition assays of A. fischeri can be used to measure for organic solvents, heavy metals, polycyclic aromatic hydrocarbons (PAH's), pesticides, and total petroleum hydrocarbons (TPH's). The bacteria’s adaptation to competitive marine environments, where they may produce unique bioactive compounds, may also position them as useful organisms for discovering novel antibiotics from marine sources. Natural transformation Natural bacterial transformation is an adaptation for transferring DNA from one individual cell to another. Natural transformation, including the uptake and incorporation of exogenous DNA into the recipient genome, has been demonstrated in A. fischeri. This process is induced by chitohexaose and is likely regulated by genes tfoX and tfoY. Natural transformation of A. fischeri facilitates rapid transfer of mutant genes across strains and provides a valuable tool for experimental genetic manipulation in this species. State microbe status In 2014, Hawaiian State Senator Glenn Wakai submitted SB3124, proposing Aliivibrio fischeri as the state microbe of Hawaii. The bill competed with a bill advocating for Flavobacterium akiainvivens to receive the same designation; ultimately, neither bill passed. In 2017, similar legislation similar to the original 2013 F. akiainvivens bill was submitted in the Hawaii House of Representatives by Isaac Choy and in the Hawaii Senate by Brian Taniguchi, but A. fischeri did not appear in this or any later proposals. List of synonyms Achromobacter fischeri (Beijerinck 1889) Bergey et al. 1930 Bacillus fischeri (Beijerinck 1889) Trevisan 1889 Bacterium phosphorescens indigenus (Eisenberg 1891) Chester 1897 Einheimischer leuchtbacillus Fischer 1888 Microspira fischeri (Beijerinck 1889) Chester 1901 Microspira marina (Russell 1892) Migula 1900 Photobacterium fischeri Beijerinck 1889 Vibrio noctiluca Weisglass and Skreb 1963	Biology and health sciences	Gram-negative bacteria	Plants
1079866	https://en.wikipedia.org/wiki/Mantle%20%28geology%29	Mantle (geology)	A mantle is a layer inside a planetary body bounded below by a core and above by a crust. Mantles are made of rock or ices, and are generally the largest and most massive layer of the planetary body. Mantles are characteristic of planetary bodies that have undergone differentiation by density. All terrestrial planets (including Earth), half of the giant planets, specifically ice giants, a number of asteroids, and some planetary moons have mantles. Examples Earth The Earth's mantle is a layer of silicate rock between the crust and the outer core. Its mass of 4.01 × 1024 kg is 67% the mass of the Earth. It has a thickness of making up about 84% of Earth's volume. It is predominantly solid, but in geological time it behaves as a viscous fluid. Partial melting of the mantle at mid-ocean ridges produces oceanic crust, and partial melting of the mantle at subduction zones produces continental crust. Other planets Mercury has a silicate mantle approximately thick, constituting only 28% of its mass. Venus's silicate mantle is approximately thick, constituting around 70% of its mass. Mars's silicate mantle is approximately thick, constituting ~74–88% of its mass, and may be represented by chassignite meteorites. Uranus and Neptune's ice mantles are approximately 30,000 km thick, composing 80% of both masses. Moons Jupiter's moons Io, Europa, and Ganymede have silicate mantles; Io's ~ silicate mantle is overlain by a volcanic crust, Ganymede's ~ thick silicate mantle is overlain by ~ of ice, and Europa's ~ km silicate mantle is overlain by ~ of ice and possibly liquid water. The silicate mantle of the Earth's moon is approximately 1300–1400 km thick, and is the source of mare basalts. The lunar mantle might be exposed in the South Pole-Aitken basin or the Crisium basin. The lunar mantle contains a seismic discontinuity at ~ depth, most likely related to a change in composition. Titan and Triton each have a mantle made of ice or other solid volatile substances. Asteroids Some of the largest asteroids have mantles; for example, Vesta has a silicate mantle similar in composition to diogenite meteorites.	Physical sciences	Geophysics	null
1080236	https://en.wikipedia.org/wiki/Dyeing	Dyeing	Dyeing is the application of dyes or pigments on textile materials such as fibers, yarns, and fabrics with the goal of achieving color with desired color fastness. Dyeing is normally done in a special solution containing dyes and particular chemical material. Dye molecules are fixed to the fiber by absorption, diffusion, or bonding with temperature and time being key controlling factors. The bond between the dye molecule and fiber may be strong or weak, depending on the dye used. Dyeing and printing are different applications; in printing, color is applied to a localized area with desired patterns. In dyeing, it is applied to the entire textile. The primary source of dye, historically, has been nature, with the dyes being extracted from plants or animals. Since the mid-19th century, however, humans have produced artificial dyes to achieve a broader range of colors and to render the dyes more stable for washing and general use. Different classes of dyes are used for different types of fiber and at different stages of the textile production process, from loose fibers through yarn and cloth to complete garments. Acrylic fibers are dyed with basic dyes, while nylon and protein fibers such as wool and silk are dyed with acid dyes, and polyester yarn is dyed with dispersed dyes. Cotton is dyed with a range of dye types, including vat dyes, and modern synthetic reactive and direct dyes. Etymology The word 'dye' (, ) comes from the Middle English , and from the Old English and . The first known use of the word 'dye' was before the 12th century. History The earliest dyed flax fibers have been found in a prehistoric cave in Georgia and dates back to 34,000 BC. More evidence of textile dyeing dates back to the Neolithic period at the large Neolithic settlement at Çatalhöyük in southern Anatolia, where traces of red dyes, possibly from ocher, an iron oxide pigment derived from clay, were found. In China, dyeing with plants, barks, and insects has been traced back more than 5,000 years. Early evidence of dyeing comes from Sindh province in Ancient India modern day Pakistan, where a piece of cotton dyed with a vegetable dye was recovered from the archaeological site at Mohenjo-daro (3rd millennium BCE). The dye used in this case was madder, which, along with other dyes such as indigo, was introduced to other regions through trade. Natural insect dyes such as Cochineal and kermes and plant-based dyes such as woad, indigo and madder were important elements of the economies of Asia and Europe until the discovery of man-made synthetic dyes in the mid-19th century. The first synthetic dye was William Perkin's mauveine in 1856, derived from coal tar. Alizarin, the red dye present in madder, was the first natural pigment to be duplicated synthetically in 1869, a development which led to the collapse of the market for naturally grown madder. The development of new, strongly colored synthetic dyes followed quickly, and by the 1870s commercial dyeing with natural dyestuffs was disappearing. An important characteristic was light-fastness - resistance to fading when exposed to sunlight using industrial techniques such as those developed by James Morton. Methods Dyeing can be applied at various stages within the textile manufacturing process; for example, fibers may be dyed before being spun into yarns, and yarns may be dyed before being woven into fabrics. Fabrics and sometimes finished garments themselves may also be dyed. The stage at which a product is dyed varies depending on its intended end use, the cost to the manufacturer, its desired appearance, and the resources available, amongst other reasons. There are specific terms to describe these dyeing methods, such as: Dope dyeing: In dope dyeing, pigments are added to the polymer solution itself before extruding the fibers. The process provides the dyed fibers with excellent fastness properties. The dope dyeing applies to synthetic fiber only. This method of dyeing is also known as solution dyeing and 'mass coloration' or 'mass colored'. It has limited color options. Fiber dyeing: In fiber dyeing, the dyeing takes place at the fiber stage before they are spun into yarn. It is also called stock dyeing. Examples are melanges and medleys. Yarn dyeing: In yarn dyeing, the yarns are dyed first before the fabric manufacturing stage. The yarn dyeing happens in hanks or in package dyeing. Package dyeing is a method where yarns are wound on perforated cones placed in a dye vessel. The dye solution is then alternatively pushed inside out and vice versa. Examples are many stripes, patterned (checks) and jacquard designed fabrics. Piece dyeing: In piece dyeing, the dyeing takes place after producing fabrics with undyed yarns. Most of the solid dyed fabrics are dyed with the piece dyeing method, and the materials are also called piece dyed. Garment dyeing: In garment dyeing, the garments are constructed of undyed, but ready-for-dyeing, fabrics. Terms for different dyed materials There are various terms used in the manufacturing and marketing industries depending on the method used to dye the substrate. For example, "stock dyed" refers to dyeing the fibers before making the yarn, "yarn dyed" refers to dyeing the yarns before producing fabrics, and "piece dyed" or "fabric dyed" refers to dyeing the yarns after they are converted into fabric. The fastness of fiber- and yarn-dyed materials is superior to that of fabrics. Objective The primary objective of the dyeing process is to achieve uniform color application in accordance with a predetermined color matching standard or reference on the substrate, which may be a fiber, yarn, or fabric, while meeting specified colour fastness requirements. Tie-dye and printing are the methods where the color is applied in a localized manner. Application Exhaust method In the exhaust method, the dye is transported to the substrate by the dye liquor's motion. The dye is adsorbed onto the fibre surface and ideally diffuses into the whole of the fibre. Water consumption in exhaust application is higher than the continuous dyeing method. There are three corresponding ways of dyeing with the exhaust method. Liquor circulating: loose stock, sliver, tow, yarn or fabric, is packed into canisters, wound onto cones or perforated beams and placed inside the dyeing vessel. In this way the liquor is pumped and revolves through the material which is stationary. Material circulating: Fabric winch dyeing and jiggers are the few forms in which material remains in motion and liquor stationary. In this the material moves through the stationary liquor. Liquor and material both in motion: Jet dyeing and softflow dyeing application methods where material and liquor both remain in motion. Continuous method In continuous method dye is transported to the substrate by passing it through the different stages but continuously. The continuous method is an innovative method where many discrete dyeing stages are combined, such as applying color, fixation and, washing off of unfixed dyes. Types of continuous dyeing are as follows Pad-steam Pad dry Thermosol Cold pad batch method is a semi-continuous dyeing process. Waterless dyeing method Waterless dyeing, also known as dry dyeing, is the newly developed and more sustainable dyeing method in which the dyes are applied to the substrate with the help of carbon dioxide or solutions that need less or no water compared to their counterparts. Selection of dyes The selection of the appropriate dyes is most important because any given dye does not apply to every type of fiber. Dyes are classified according to many parameters, such as chemical structure, affinity, application method, desired colour fastness i.e. resistance to washing, rubbing, and light. The properties may vary with different dyes. The selection of dye depends on the objective in dyeing and affinity (to which material is to be dyed). Fastness of color largely depends upon the molecular size of the dyes and the solubility. Larger molecular size serves better washing fastness results. Indigo dyes have a poor wash and rubbing fastness on denim (cotton), so they are used to produce washed-down effects on fabrics. In contrast, vat or reactive dyes are applied to cotton to achieve excellent washing fastness. The next important criterion for selecting dyes is the assessment of hazards to human health and the environment. There are many dyes especially disperse dyes that may cause allergic reactions to some individuals, and the negative impact on the environment. There are national and international standards and regulations which need to comply. Direct application The term "direct dye application" stems from some dyestuff having to be either fermented as in the case of some natural dye or chemically reduced as in the case of synthetic vat and sulfur dyes before being applied. This renders the dye soluble so that it can be absorbed by the fiber since the insoluble dye has very little substantivity to the fiber. Direct dyes, a class of dyes largely for dyeing cotton, are water-soluble and can be applied directly to the fiber from an aqueous solution. Most other classes of synthetic dye, other than vat and surface dyes, are also applied in this way. The term may also be applied to dyeing without the use of mordants to fix the dye once it is applied. Mordants were often required to alter the hue and intensity of natural dyes and improve color fastness. Chromium salts were until recently extensively used in dyeing wool with synthetic mordant dyes. These were used for economical high color fastness dark shades such as black and navy. Environmental concerns have now restricted their use, and they have been replaced with reactive and metal complex dyes that do not require mordant. Yarn dyeing There are many forms of yarn dyeing. Common forms are the package form and the hanks form. Cotton yarns are mostly dyed at package form, and acrylic or wool yarn are dyed at hank form. In the continuous filament industry, polyester or polyamide yarns are always dyed at package form, while viscose rayon yarns are partly dyed at hank form because of technology. The common dyeing process of cotton yarn with reactive dyes at package form is as follows: The raw yarn is wound on a spring tube to achieve a package suitable for dye penetration. These softened packages are loaded on a dyeing carrier's spindle one on another. The packages are pressed up to a desired height to achieve suitable density of packing. The carrier is loaded on the dyeing machine and the yarn is dyed. After dyeing, the packages are unloaded from the carrier into a trolley. Now the trolley is taken to hydro extractor where water is removed. The packages are hydro extracted to remove the maximum amount of water leaving the desired color into raw yarn. The packages are then dried to achieve the final dyed package. After this process, the dyed yarn packages are packed and delivered. Space dyeing Space dyeing is a technique of localized color application that produces a unique multicolored effect. History of garment dyeing Garment dyeing is the process of dyeing fully fashioned garments subsequent to manufacturing, as opposed to the conventional method of manufacturing garments from pre-dyed fabrics. Up until the mid-1970s the method was rarely used for commercial clothing production. It was used domestically, to overdye old, worn and faded clothes, and also by resellers of used or surplus military clothing. The first notable industrial use of the technique was made by Benetton, which garment dyed its Shetland wool knitwear. Complex garment dyeing In the mid-1970s the Bologna clothing designer Massimo Osti began experimenting with the garment dyeing technique. His experimentation over the next decade, led to the pioneering of not just the industrial use of traditional garment dyeing (dyeing simple cotton or wool garments) but, more importantly, the technique of “complex garment dyeing” which involved dyeing fully fashioned garments which had been constructed from multiple fabric or fiber types (e.g. a jacket made from both nylon and cotton, or linen, nylon and polyurethane coated cotton) in the same bath. Up until its development by Osti (for his clothing brand C.P. Company), this technique had never been successfully industrially applied in any context. The complexity lay in developing both a practical and chemical understanding of how each fabric responded differently to the dye, how much it would shrink, how much color it would absorb, developing entirely new forms of quality control to verify possible defects in fabric before dyeing etc. Beyond the industrial advantages of the technique (purchasing fabric in one color, white or natural, you may produce as many colors as you wish etc.), the artistic advantages of the technique were considerable and in many ways paved the way for the creation of the clothing style today known as Italian Sportswear. These advantages included the way in which different fibers absorbed the dye's color allowed for the creation of incredibly nuanced differences in color tones and a harmony that is impossible to achieve any other way the garment dyeing process naturally gave the fabric a “worn-in” hand allowing for the development of the casual and relaxed version of the classic menswear look which characterizes Italian sportswear the fact that each fabric and fiber type responds differently to the dye also produces a “deconstructed” effect, whereby the consumer's attention is drawn to the construction techniques of the jacket. For example: a more densely woven fabric absorbs the color less intensely than a more open weave, the polyester stitching used for a cotton garment does not absorb any dye color, producing a contrast color stitch etc. The disadvantages included: a relatively high failure rate for garments (between 5–10%) the difficulty in achieving a very tailored look due to difficulties in precisely calculating shrinkage rates high research and prototyping costs in order to understand how fabrics will behave in production Today, whilst garment dyeing is a diffusely employed as an industrial technique around the globe, predominantly in the production of vintage style cotton garments and by fast fashion suppliers, complex garment dyeing is still practiced almost exclusively in Italy, by a handful of premium brands and suppliers who remain committed to the art. Related terms There are several terms associated with the process of dyeing: Affinity Affinity refers to the chemical attraction between two elements or substances, leading to their inclination to unite or combine, as observed between fiber and dyestuff. Bleeding Materials that exhibit bleeding tendencies may lead to the staining of white or light-colored fabrics in contact with them while in a wet state. The phenomenon of color fading from a fabric or yarn upon immersion in water, solvent, or a comparable liquid medium, arises due to inadequate dyeing or the utilization of inferior quality dyes. Staining Fabric can experience undesired color absorption, resulting in staining, when exposed to water, dry-cleaning solvent, or similar liquids containing unintended dyestuffs or coloring materials. Additionally, direct contact with other dyed materials may cause color transfer through bleeding or sublimation. Stripping Stripping is a method used to partially or entirely remove color from dyed textile materials. It can also be utilized as a reprocessing technique to correct imperfect dyeing.	Technology	Other techniques	null
1081235	https://en.wikipedia.org/wiki/Lava%20dome	Lava dome	In volcanology, a lava dome is a circular, mound-shaped protrusion resulting from the slow extrusion of viscous lava from a volcano. Dome-building eruptions are common, particularly in convergent plate boundary settings. Around 6% of eruptions on Earth form lava domes. The geochemistry of lava domes can vary from basalt (e.g. Semeru, 1946) to rhyolite (e.g. Chaiten, 2010) although the majority are of intermediate composition (such as Santiaguito, dacite-andesite, present day) The characteristic dome shape is attributed to high viscosity that prevents the lava from flowing very far. This high viscosity can be obtained in two ways: by high levels of silica in the magma, or by degassing of fluid magma. Since viscous basaltic and andesitic domes weather fast and easily break apart by further input of fluid lava, most of the preserved domes have high silica content and consist of rhyolite or dacite. Existence of lava domes has been suggested for some domed structures on the Moon, Venus, and Mars, e.g. the Martian surface in the western part of Arcadia Planitia and within Terra Sirenum. Dome dynamics Lava domes evolve unpredictably, due to non-linear dynamics caused by crystallization and outgassing of the highly viscous lava in the dome's conduit. Domes undergo various processes such as growth, collapse, solidification and erosion. Lava domes grow by endogenic dome growth or exogenic dome growth. The former implies the enlargement of a lava dome due to the influx of magma into the dome interior, and the latter refers to discrete lobes of lava emplaced upon the surface of the dome. It is the high viscosity of the lava that prevents it from flowing far from the vent from which it extrudes, creating a dome-like shape of sticky lava that then cools slowly in-situ. Spines and lava flows are common extrusive products of lava domes. Domes may reach heights of several hundred meters, and can grow slowly and steadily for months (e.g. Unzen volcano), years (e.g. Soufrière Hills volcano), or even centuries (e.g. Mount Merapi volcano). The sides of these structures are composed of unstable rock debris. Due to the intermittent buildup of gas pressure, erupting domes can often experience episodes of explosive eruption over time. If part of a lava dome collapses and exposes pressurized magma, pyroclastic flows can be produced. Other hazards associated with lava domes are the destruction of property from lava flows, forest fires, and lahars triggered from re-mobilization of loose ash and debris. Lava domes are one of the principal structural features of many stratovolcanoes worldwide. Lava domes are prone to unusually dangerous explosions since they can contain rhyolitic silica-rich lava. Characteristics of lava dome eruptions include shallow, long-period and hybrid seismicity, which is attributed to excess fluid pressures in the contributing vent chamber. Other characteristics of lava domes include their hemispherical dome shape, cycles of dome growth over long periods, and sudden onsets of violent explosive activity. The average rate of dome growth may be used as a rough indicator of magma supply, but it shows no systematic relationship to the timing or characteristics of lava dome explosions. Gravitational collapse of a lava dome can produce a block and ash flow. Related landforms Cryptodomes A cryptodome (from the Greek , , "hidden, secret") is a dome-shaped structure created by accumulation of viscous magma at a shallow depth. One example of a cryptodome was in the May 1980 eruption of Mount St. Helens, where the explosive eruption began after a landslide caused the side of the volcano to collapse, leading to explosive decompression of the subterranean cryptodome. Lava spine/Lava spire A lava spine or lava spire is a growth that can form on the top of a lava dome. A lava spine can increase the instability of the underlying lava dome. A recent example of a lava spine is the spine formed in 1997 at the Soufrière Hills Volcano on Montserrat. Lava coulées Coulées (or coulees) are lava domes that have experienced some flow away from their original position, thus resembling both lava domes and lava flows. The world's largest known dacite flow is the Chao dacite dome complex, a huge coulée flow-dome between two volcanoes in northern Chile. This flow is over long, has obvious flow features like pressure ridges, and a flow front tall (the dark scalloped line at lower left). There is another prominent coulée flow on the flank of Llullaillaco volcano, in Argentina, and other examples in the Andes. Examples of lava domes	Physical sciences	Volcanology	Earth science
1082009	https://en.wikipedia.org/wiki/Sand%20casting	Sand casting	Sand casting, also known as sand molded casting, is a metal casting process characterized by using sand—known as casting sand—as the mold material. The term "sand casting" can also refer to an object produced via the sand casting process. Sand castings are produced in specialized factories called foundries. In 2003, over 60% of all metal castings were produced via sand casting. Molds made of sand are relatively cheap, and sufficiently refractory even for steel foundry use. In addition to the sand, a suitable bonding agent (usually clay) is mixed or occurs with the sand. The mixture is moistened, typically with water, but sometimes with other substances, to develop the strength and plasticity of the clay and to make the aggregate suitable for molding. The sand is typically contained in a system of frames or mold boxes known as a flask. The mold cavities and gate system are created by compacting the sand around models called patterns, by carving directly into the sand, or via 3D printing. Basic process There are five steps in this process: Place a pattern in sand to create a mold. Incorporate the pattern and sand in a gating system. Remove the pattern. Fill the mold cavity with molten metal. Allow the metal to cool. Break away the sand mold and remove the casting. Components Patterns From the design, provided by a designer, a skilled pattern maker builds a pattern of the object to be produced, using wood, metal, or a plastic such as expanded polystyrene. Sand can be ground, swept or strickled into shape. The metal to be cast will contract during solidification, and this may be non-uniform due to uneven cooling. Therefore, the pattern must be slightly larger than the finished product, a difference known as contraction allowance. Different scaled rules are used for different metals, because each metal and alloy contracts by an amount distinct from all others. Patterns also have core prints that create registers within the molds into which are placed sand cores. Such cores, sometimes reinforced by wires, are used to create under-cut profiles and cavities which cannot be molded with the cope and drag, such as the interior passages of valves or cooling passages in engine blocks. Paths for the entrance of metal into the mold cavity constitute the runner system and include the sprue, various feeders which maintain a good metal 'feed', and in-gates which attach the runner system to the casting cavity. Gas and steam generated during casting exit through the permeable sand or via risers, which are added either in the pattern itself, or as separate pieces. Tools In addition to patterns, the sand molder could also use tools to create the holes. Molding box and materials A multi-part molding box (known as a casting flask, the top and bottom halves of which are known respectively as the cope and drag) is prepared to receive the pattern. Molding boxes are made in segments that may be latched to each other and to end closures. For a simple object—flat on one side—the lower portion of the box, closed at the bottom, will be filled with a molding sand. The sand is packed in through a vibratory process called ramming, and in this case, periodically screeded level. The surface of the sand may then be stabilized with a sizing compound. The pattern is placed on the sand and another molding box segment is added. Additional sand is rammed over and around the pattern. Finally a cover is placed on the box and it is turned and unlatched, so that the halves of the mold may be parted and the pattern with its sprue and vent patterns removed. Additional sizing may be added and any defects introduced by the removal of the pattern are corrected. The box is closed again. This forms a "green" mold which must be dried to receive the hot metal. If the mold is not sufficiently dried a steam explosion can occur that can throw molten metal about. In some cases, the sand may be oiled instead of moistened, which makes casting possible without waiting for the sand to dry. Sand may also be bonded by chemical binders, such as furane resins or amine-hardened resins. Additive manufacturing (AM) can be used in the sand mold preparation, so that instead of the sand mold being formed via packing sand around a pattern, it is 3D-printed. This can reduce lead times for casting by obviating patternmaking. Besides replacing older methods, additive can also complement them in hybrid models, such as making a variety of AM-printed cores for a cavity derived from a traditional pattern. Chills To control the solidification structure of the metal, it is possible to place metal plates, chills, in the mold. The associated rapid local cooling will form a finer-grained structure and may form a somewhat harder metal at these locations. In ferrous castings, the effect is similar to quenching metals in forge work. The inner diameter of an engine cylinder is made hard by a chilling core. In other metals, chills may be used to promote directional solidification of the casting. In controlling the way a casting freezes, it is possible to prevent internal voids or porosity inside castings. Cores Cores are apparatus used to generate hollow cavities or internal features which cannot be formed using pattern alone in moulding, cores are usually made using sand, but some processes also use permanent cores made of metal. To produce cavities within the casting—such as for liquid cooling in engine blocks and cylinder heads—negative forms are used to produce cores. Usually sand-molded, cores are inserted into the casting box after removal of the pattern. Whenever possible, designs are made that avoid the use of cores, due to the additional set-up time, mass and thus greater cost. With a completed mold at the appropriate moisture content, the box containing the sand mold is then positioned for filling with molten metal—typically iron, steel, bronze, brass, aluminium, magnesium alloys, or various pot metal alloys, which often include lead, tin, and zinc. After being filled with liquid metal the box is set aside until the metal is sufficiently cool to be strong. The sand is then removed, revealing a rough casting that, in the case of iron or steel, may still be glowing red. In the case of metals that are significantly heavier than the casting sand, such as iron or lead, the casting flask is often covered with a heavy plate to prevent a problem known as floating the mold. Floating the mold occurs when the pressure of the metal pushes the sand above the mold cavity out of shape, causing the casting to fail. After casting, the cores are broken up by rods or shot and removed from the casting. The metal from the sprue and risers is cut from the rough casting. Various heat treatments may be applied to relieve stresses from the initial cooling and to add hardness—in the case of steel or iron, by quenching in water or oil. The casting may be further strengthened by surface compression treatment—like shot peening—that adds resistance to tensile cracking and smooths the rough surface. And when high precision is required, various machining operations (such as milling or boring) are made to finish critical areas of the casting. Examples of this would include the boring of cylinders and milling of the deck on a cast engine block. Design requirements The part to be made and its pattern must be designed to accommodate each stage of the process, as it must be possible to remove the pattern without disturbing the molding sand and to have proper locations to receive and position the cores. A slight taper, known as draft, must be used on surfaces perpendicular to the parting line, in order to be able to remove the pattern from the mold. This requirement also applies to cores, as they must be removed from the core box in which they are formed. The sprue and risers must be arranged to allow a proper flow of metal and gasses within the mold in order to avoid an incomplete casting. Should a piece of core or mold become dislodged it may be embedded in the final casting, forming a sand pit, which may render the casting unusable. Gas pockets can cause internal voids. These may be immediately visible or may only be revealed after extensive machining has been performed. For critical applications, or where the cost of wasted effort is a factor, non-destructive testing methods may be applied before further work is performed. Processes In general, we can distinguish between two methods of sand casting; the first one using green sand and the second being the air set method. Green sand These castings are made using sand molds formed from "wet" sand which contains water and organic bonding compounds, typically referred to as clay. The name "green sand" comes from the fact that the sand mold is not "set", it is still in the "green" or uncured state even when the metal is poured in the mould. Green sand is not green in color, but "green" in the sense that it is used in a wet state (akin to green wood). Contrary to what the name suggests, "green sand" is not a type of sand on its own (that is, not greensand in the geologic sense), but is rather a mixture of: silica sand (SiO2), chromite sand (FeCr2O4), or zircon sand (ZrSiO4), 75 to 85%, sometimes with a proportion of olivine, staurolite, or graphite. bentonite (clay), 5 to 11% water, 2 to 4% inert sludge 3 to 5% anthracite (0 to 1%) There are many recipes for the proportion of clay, but they all strike different balances between moldability, surface finish, and ability of the hot molten metal to degas. Coal, typically referred to in foundries as sea-coal, which is present at a ratio of less than 5%, partially combusts in the presence of the molten metal, leading to offgassing of organic vapors. Green sand casting for non-ferrous metals does not use coal additives, since the CO created does not prevent oxidation. Green sand for aluminum typically uses olivine sand (a mixture of the minerals forsterite and fayalite, which is made by crushing dunite rock). The choice of sand has a lot to do with the temperature at which the metal is poured. At the temperatures that copper and iron are poured, the clay is inactivated by the heat, in that the montmorillonite is converted to illite, which is a non-expanding clay. Most foundries do not have the very expensive equipment to remove the burned out clay and substitute new clay, so instead, those that pour iron typically work with silica sand that is inexpensive compared to the other sands. As the clay is burned out, newly mixed sand is added and some of the old sand is discarded or recycled into other uses. Silica is the least desirable of the sands, since metamorphic grains of silica sand have a tendency to explode to form sub-micron sized particles when thermally shocked during pouring of the molds. These particles enter the air of the work area and can lead to silicosis in the workers. Iron foundries expend considerable effort on aggressive dust collection to capture this fine silica. Various types of respiratory-protective equipment are also used in foundries. The sand also has the dimensional instability associated with the conversion of quartz from alpha quartz to beta quartz at 680 °C (1250 °F). Often, combustible additives such as wood flour are added to create spaces for the grains to expand without deforming the mold. Olivine, chromite, etc. are therefore used because they do not have a phase transition that causes rapid expansion of the grains. Olivine and chromite also offer greater density, which cools the metal faster, thereby producing finer grain structures in the metal. Since they are not metamorphic minerals, they do not have the polycrystals found in silica, and subsequently they do not form hazardous sub-micron sized particles. "Air set" method The air set method uses dry sand bonded with materials other than clay, using a fast curing adhesive. The latter may also be referred to as no bake mold casting. When these are used, they are collectively called "air set" sand castings to distinguish them from "green sand" castings. Two types of molding sand are natural bonded (bank sand) and synthetic (lake sand); the latter is generally preferred due to its more consistent composition. With both methods, the sand mixture is packed around a pattern, forming a mold cavity. If necessary, a temporary plug is placed in the sand and touching the pattern in order to later form a channel into which the casting fluid can be poured. Air-set molds are often formed with the help of a casting flask having a top and bottom part, termed the cope and drag. The sand mixture is tamped down as it is added around the pattern, and the final mold assembly is sometimes vibrated to compact the sand and fill any unwanted voids in the mold. Then the pattern is removed along with the channel plug, leaving the mold cavity. The casting liquid (typically molten metal) is then poured into the mold cavity. After the metal has solidified and cooled, the casting is separated from the sand mold. There is typically no mold release agent, and the mold is generally destroyed in the removal process. The accuracy of the casting is limited by the type of sand and the molding process. Sand castings made from coarse green sand impart a rough texture to the surface, and this makes them easy to identify. Castings made from fine green sand can shine as cast but are limited by the depth to width ratio of pockets in the pattern. Air-set molds can produce castings with smoother surfaces than coarse green sand but this method is primarily chosen when deep narrow pockets in the pattern are necessary, due to the expense of the plastic used in the process. Air-set castings can typically be easily identified by the burnt color on the surface. The castings are typically shot blasted to remove that burnt color. Surfaces can also be later ground and polished, for example when making a large bell. After molding, the casting is covered with a residue of oxides, silicates and other compounds. This residue can be removed by various means, such as grinding, or shot blasting. During casting, some of the components of the sand mixture are lost in the thermal casting process. Green sand can be reused after adjusting its composition to replenish the lost moisture and additives. The pattern itself can be reused indefinitely to produce new sand molds. The sand molding process has been used for many centuries to produce castings manually. Since 1950, partially automated casting processes have been developed for production lines. Cold box Cold box uses organic and inorganic binders that strengthen the mold by chemically adhering to the sand. This type of mold gets its name from not being baked in an oven like other sand mold types. This type of mold is more accurate dimensionally than green-sand molds but is more expensive. Thus it is used only in applications that necessitate it. No-bake molds No-bake molds are expendable sand molds, similar to typical sand molds, except they also contain a quick-setting liquid resin and catalyst. Rather than being rammed, the molding sand is poured into the flask and held until the resin solidifies, which occurs at room temperature. This type of molding also produces a better surface finish than other types of sand molds. Because no heat is involved it is called a cold-setting process. Common flask materials that are used are wood, metal, and plastic. Common metals cast into no-bake molds are brass, iron (ferrous), and aluminum alloys. Vacuum molding Vacuum molding (V-process) is a variation of the sand casting process for most ferrous and non-ferrous metals, in which unbonded sand is held in the flask with a vacuum. The pattern is specially vented so that a vacuum can be pulled through it. A heat-softened thin sheet () of plastic film is draped over the pattern and a vacuum is drawn (). A special vacuum forming flask is placed over the plastic pattern and is filled with a free-flowing sand. The sand is vibrated to compact the sand and a sprue and pouring cup are formed in the cope. Another sheet of plastic is placed over the top of the sand in the flask and a vacuum is drawn through the special flask; this hardens and strengthens the unbonded sand. The vacuum is then released on the pattern and the cope is removed. The drag is made in the same way (without the sprue and pouring cup). Any cores are set in place and the mold is closed. The molten metal is poured while the cope and drag are still under a vacuum, because the plastic vaporizes but the vacuum keeps the shape of the sand while the metal solidifies. When the metal has solidified, the vacuum is turned off and the sand runs out freely, releasing the casting. The V-process is known for not requiring a draft because the plastic film has a certain degree of lubricity and it expands slightly when the vacuum is drawn in the flask. The process has high dimensional accuracy, with a tolerance of ±0.010 in for the first inch and ±0.002 in/in thereafter. Cross-sections as small as are possible. The surface finish is very good, usually between 150 and 125 rms. Other advantages include no moisture related defects, no cost for binders, excellent sand permeability, and no toxic fumes from burning the binders. Finally, the pattern does not wear out because the sand does not touch it. The main disadvantage is that the process is slower than traditional sand casting so it is only suitable for low to medium production volumes; approximately 10 to 15,000 pieces a year. However, this makes it perfect for prototype work, because the pattern can be easily modified as it is made from plastic. Fast mold making processes With the fast development of the car and machine building industry the casting consuming areas called for steady higher productivity. The basic process stages of the mechanical molding and casting process are similar to those described under the manual sand casting process. The technical and mental development however was so rapid and profound that the character of the sand casting process changed radically. Mechanized sand molding The first mechanized molding lines consisted of sand slingers and/or jolt-squeeze devices that compacted the sand in the flasks. Subsequent mold handling was mechanical using cranes, hoists and straps. After core setting the copes and drags were coupled using guide pins and clamped for closer accuracy. The molds were manually pushed off on a roller conveyor for casting and cooling. Automatic high pressure sand molding lines Increasing quality requirements made it necessary to increase the mold stability by applying steadily higher squeeze pressure and modern compaction methods for the sand in the flasks. In early fifties the high pressure molding was developed and applied in mechanical and later automatic flask lines. The first lines were using jolting and vibrations to pre-compact the sand in the flasks and compressed air powered pistons to compact the molds. Horizontal sand flask molding In the first automatic horizontal flask lines the sand was shot or slung down on the pattern in a flask and squeezed with hydraulic pressure of up to 140 bars. The subsequent mold handling including turn-over, assembling, pushing-out on a conveyor were accomplished either manually or automatically. In the late fifties hydraulically powered pistons or multi-piston systems were used for the sand compaction in the flasks. This method produced much more stable and accurate molds than it was possible manually or pneumatically. In the late sixties mold compaction by fast air pressure or gas pressure drop over the pre-compacted sand mold was developed (sand-impulse and gas-impact). The general working principle for most of the horizontal flask line systems is shown on the sketch below. Today there are many manufacturers of the automatic horizontal flask molding lines. The major disadvantages of these systems is high spare parts consumption due to multitude of movable parts, need of storing, transporting and maintaining the flasks and productivity limited to approximately 90–120 molds per hour. Vertical sand flaskless molding In 1962, Dansk Industri Syndikat A/S (DISA-DISAMATIC) invented a flask-less molding process by using vertically parted and poured molds. The first line could produce up to 240 complete sand molds per hour. Today molding lines can achieve a molding rate of 550 sand molds per hour and requires only one monitoring operator. Maximum mismatch of two mold halves is . Although very fast, vertically parted molds are not typically used by jobbing foundries due to the specialized tooling needed to run on these machines. Cores need to be set with a core mask as opposed to by hand and must hang in the mold as opposed to being set on parting surface. Matchplate sand molding The principle of the matchplate, meaning pattern plates with two patterns on each side of the same plate, was developed and patented in 1910, fostering the perspectives for future sand molding improvements. However, first in the early sixties the American company Hunter Automated Machinery Corporation launched its first automatic flaskless, horizontal molding line applying the matchplate technology. The method alike to the DISA's (DISAMATIC) vertical molding is flaskless, however horizontal. The matchplate molding technology is today used widely. Its great advantage is inexpensive pattern tooling, easiness of changing the molding tooling, thus suitability for manufacturing castings in short series so typical for the jobbing foundries. Modern matchplate molding machine is capable of high molding quality, less casting shift due to machine-mold mismatch (in some cases less than ), consistently stable molds for less grinding and improved parting line definition. In addition, the machines are enclosed for a cleaner, quieter working environment with reduced operator exposure to safety risks or service-related problems. Safety standards With automated mold manufacturing came additional workplace safety requirements. Different voluntary technical standards apply depending on the geopolitical jurisdiction where the machinery is to be used. Canada Canada does not have a machine-specific voluntary technical standard for sand-mold making machinery. This type of machinery is covered by: Safeguarding of machinery, CSA Z432. Canadian Standards Association. 2016. In addition, the electrical safety requirements are covered by: Industrial Electrical Machinery, CSA C22.2 No. 301. 2016. European Union The primary standard for sand-mold manufacturing equipment in the EU is: Safety requirements for foundry moulding and coremaking machinery and plant associated equipment, EN 710. European Committee for Standardization (CEN). EN 710 will need to be used in conjunction with EN 60204-1 for electrical safety, and EN ISO 13849-1 and EN ISO 13849-2 or EN 62061 for functional safety. Additional type C standards may also be necessary for conveyors, robotics or other equipment that may be needed to support the operation of the mold-making equipment. United States There is no machine-specific standard for sand-mold manufacturing equipment. The ANSI B11 family of standards includes some generic machine-tool standards that could be applied to this type of machinery, including: Safety of Machinery, ANSI B11.0. American National Standards Institute (ANSI). 2020. Performance Requirements for Risk Reduction Measures: Safeguarding and other Means of Reducing Risk, ANSI B11.19. American National Standards Institute (ANSI). 2019. Safety Requirements for the Integration of Machinery into a System, ANSI B11.20. American National Standards Institute (ANSI). 2017. Safety Requirements for Transfer Machines, ANSI B11.24. American National Standards Institute (ANSI). 2002 (R2020). Functional Safety for Equipment (Electrical/Fluid Power Control Systems) General Principles for the Design of Safety Control Systems Using ISO 13849-1, ANSI B11.26. American National Standards Institute (ANSI). 2018. Sound Level Measurement Guidelines, ANSI B11.TR5. American National Standards Institute (ANSI). 2006 (R2017). Mold materials There are four main components for making a sand casting mold: base sand, a binder, additives, and a parting compound. Molding sands Molding sands, also known as foundry sands, are defined by eight characteristics: refractoriness, chemical inertness, permeability, surface finish, cohesiveness, flowability, collapsibility, and availability/cost. Refractoriness — This refers to the sand's ability to withstand the temperature of the liquid metal being cast without breaking down. For example, some sands only need to withstand if casting aluminum alloys, whereas steel needs a sand that will withstand . Sand with too low refractoriness will melt and fuse to the casting. Chemical inertness — The sand must not react with the metal being cast. This is especially important with highly reactive metals, such as magnesium and titanium. Permeability — This refers to the sand's ability to exhaust gases. This is important because during the pouring process many gases are produced, such as hydrogen, nitrogen, carbon dioxide, and steam, which must leave the mold otherwise casting defects, such as blow holes and gas holes, occur in the casting. Note that for each cubic centimeter (cc) of water added to the mold 1600 cc of steam is produced. Surface finish — The size and shape of the sand particles defines the best surface finish achievable, with finer particles producing a better finish. However, as the particles become finer (and surface finish improves) the permeability becomes worse. Cohesiveness (or bond) — This is the ability of the sand to retain a given shape after the pattern is removed. Flowability – The ability for the sand to flow into intricate details and tight corners without special processes or equipment. Collapsibility — This is the ability of the sand to be easily stripped off the casting after it has solidified. Sands with poor collapsibility will adhere strongly to the casting. When casting metals that contract a lot during cooling or with long freezing temperature ranges a sand with poor collapsibility will cause cracking and hot tears in the casting. Special additives can be used to improve collapsibility. Availability/cost — The availability and cost of the sand is very important because for every ton of metal poured, three to six tons of sand is required. Although sand can be screened and reused, the particles eventually become too fine and require periodic replacement with fresh sand. In large castings it is economical to use two different sands, because the majority of the sand will not be in contact with the casting, so it does not need any special properties. The sand that is in contact with the casting is called facing sand, and is designed for the casting on hand. This sand will be built up around the pattern to a thickness of . The sand that fills in around the facing sand is called backing sand. This sand is simply silica sand with only a small amount of binder and no special additives. Types of base sands Base sand is the type used to make the mold or core without any binder. Because it does not have a binder it will not bond together and is not usable in this state. Silica sand Silica (SiO2) sand is the sand found on a beach and is also the most commonly used sand. It is either made by crushing sandstone or taken from natural occurring locations, such as beaches and river beds. The fusion point of pure silica is , however the sands used have a lower melting point due to impurities. For high melting point casting, such as steels, a minimum of 98% pure silica sand must be used; however for lower melting point metals, such as cast iron and non-ferrous metals, a lower purity sand can be used (between 94 and 98% pure). Silica sand is the most commonly used sand because of its great abundance, and, thus, low cost (therein being its greatest advantage). Its disadvantages are high thermal expansion, which can cause casting defects with high melting point metals, and low thermal conductivity, which can lead to unsound casting. It also cannot be used with certain basic metals because it will chemically interact with the metal, forming surface defects. Finally, it releases silica particulates during the pour, risking silicosis in foundry workers. Olivine sand Olivine is a mixture of orthosilicates of iron and magnesium from the mineral dunite. Its main advantage is that it is free from silica, therefore it can be used with basic metals, such as manganese steels. Other advantages include a low thermal expansion, high thermal conductivity, and high fusion point. Finally, it is safer to use than silica, therefore it is popular in Europe. Chromite sand Chromite sand is a solid solution of spinels. Its advantages are a low percentage of silica, a very high fusion point (), and a very high thermal conductivity. Its disadvantage is its costliness, therefore it is only used with expensive alloy steel casting and to make cores. Zircon sand Zircon sand is a compound of approximately two-thirds zirconium oxide (ZrO2) and one-third silica. It has the highest fusion point of all the base sands at , a very low thermal expansion, and a high thermal conductivity. Because of these good properties it is commonly used when casting alloy steels and other expensive alloys. It is also used as a mold wash (a coating applied to the molding cavity) to improve surface finish. However, it is expensive and not readily available. Chamotte sand Chamotte is made by calcining fire clay (Al2O3-SiO2) above . Its fusion point is and has low thermal expansion. It is the second cheapest sand, however it is still twice as expensive as silica. Its disadvantages are very coarse grains, which result in a poor surface finish, and it is limited to dry sand molding. Mold washes are used to overcome the surface finish problems. This sand is usually used when casting large steel workpieces. Binders Binders are added to a base sand to bond the sand particles together (i.e. it is the glue that holds the mold together). Clay and water A mixture of clay and water is the most commonly used binder. There are two types of clay commonly used: bentonite and kaolinite, with the former being the most common. Oil Oils, such as linseed oil, other vegetable oils and marine oils, used to be used as a binder, however due to their increasing cost, they have been mostly phased out. The oil also required careful baking at to cure (if overheated, the oil becomes brittle, wasting the mold). Resin Resin binders are natural or synthetic high melting point gums. The two common types used are urea formaldehyde (UF) and phenol formaldehyde (PF) resins. PF resins have a higher heat resistance than UF resins and cost less. There are also cold-set resins, which use a catalyst instead of a heat to cure the binder. Resin binders are quite popular because different properties can be achieved by mixing with various additives. Other advantages include good collapsibility, low gassing, and they leave a good surface finish on the casting. MDI (methylene diphenyl diisocyanate) is also a commonly used binder resin in the foundry core process. Sodium silicate Water glass ( sodium silicate [Na2SiO3 or (Na2O)(SiO2)] ) is a high strength binder used with silica molding sand both for cores and molds. To cure a mixture of finely ground sand (e.g. by using a sand muller) and 3 to 4% of sodium silicate the binder, carbon dioxide (CO2) gas is used. The mixture is exposed to the gas at ambient temperature reacting as following: {Na2O(SiO2)} + CO2 <=> {Na2CO3} + {2SiO2} + Heat The advantage to this binder is that it can be used at room temperature and is fast. The disadvantage is that its high strength leads to shakeout difficulties and possibly hot tears (probably due to quartz inversion) in the casting. The mixed sodium silicate and sand may also be heated by a heat gun to achieve better rigideness. Additives Additives are added to the molding components to improve: surface finish, dry strength, refractoriness, and "cushioning properties". Up to 5% of reducing agents, such as coal powder, pitch, creosote, and fuel oil, may be added to the molding material to prevent wetting (prevention of liquid metal sticking to sand particles, thus leaving them on the casting surface), improve surface finish, decrease metal penetration, and burn-on defects. These additives achieve this by creating gases at the surface of the mold cavity, which prevent the liquid metal from adhering to the sand. Reducing agents are not used with steel casting, because they can carburize the metal during casting. Up to 3% of "cushioning material", such as wood flour, sawdust, powdered husks, peat, and straw, can be added to reduce scabbing, hot tear, and hot crack casting defects when casting high temperature metals. These materials are beneficial because burn-off when the metal is poured creates tiny voids in the mold, allowing the sand particles to expand. They also increase collapsibility and reduce shakeout time. Up to 2% of cereal binders, such as dextrin, starch, sulphite lye, and molasses, can be used to increase dry strength (the strength of the mold after curing) and improve surface finish. Cereal binders also improve collapsibility and reduce shakeout time because they burn off when the metal is poured. The disadvantage to cereal binders is that they are expensive. Up to 2% of iron oxide powder can be used to prevent mold cracking and metal penetration, essentially improving refractoriness. Silica flour (fine silica) and zircon flour also improve refractoriness, especially in ferrous castings. The disadvantages to these additives is that they greatly reduce permeability. Parting compounds To get the pattern out of the mold, prior to casting, a parting compound is applied to the pattern to ease removal. They can be a liquid or a fine powder (particle diameters between ). Common powders include talc, graphite, and dry silica; common liquids include mineral oil and water-based silicon solutions. The latter are more commonly used with metal and large wooden patterns. History Clay molds were used in ancient China since the Shang dynasty ( to 1046 BC). The famous Houmuwu ding (c. 1300 BC) was made using clay molding. The Assyrian king Sennacherib (704–681 BC) cast massive bronzes of up to 30 tonnes, and claims to have been the first to have used clay molds rather than the "lost-wax" method: Whereas in former times the kings my forefathers had created bronze statues imitating real-life forms to put on display inside their temples, but in their method of work they had exhausted all the craftsmen, for lack of skill and failure to understand the principles they needed so much oil, wax and tallow for the work that they caused a shortage in their own countries—I, Sennacherib, leader of all princes, knowledgeable in all kinds of work, took much advice and deep thought over doing that work. Great pillars of bronze, colossal striding lions, such as no previous king had ever constructed before me, with the technical skill that Ninushki brought to perfection in me, and at the prompting of my intelligence and the desire of my heart I invented a technique for bronze and made it skillfully. I created clay moulds as if by divine intelligence....twelve fierce lion-colossi together with twelve mighty bull-colossi which were perfect castings... I poured copper into them over and over again; I made the castings as skillfully as if they had only weighed half a shekel each In 1206, Ismail al-Jazari first described the casting of metals in closed mold boxes with sand. Sand casting molding method was recorded by Vannoccio Biringuccio in his book published around 1540. In 1924, the Ford Motor Company set a record by producing 1 million cars, in the process consuming one-third of the total casting production in the U.S. As the automobile industry grew the need for increased casting efficiency grew. The increasing demand for castings in the growing car and machine building industry during and after World War I and World War II, stimulated new inventions in mechanization and later automation of the sand casting process technology. There was not one bottleneck to faster casting production but rather several. Improvements were made in molding speed, molding sand preparation, sand mixing, core manufacturing processes, and the slow metal melting rate in cupola furnaces. In 1912, the sand slinger was invented by the American company Beardsley & Piper. In 1912, the first sand mixer with individually mounted revolving plows was marketed by the Simpson Company. In 1915, the first experiments started with bentonite clay instead of simple fire clay as the bonding additive to the molding sand. This increased tremendously the green and dry strength of the molds. In 1918, the first fully automated foundry for fabricating hand grenades for the U.S. Army went into production. In the 1930s the first high-frequency coreless electric furnace was installed in the U.S. In 1943, ductile iron was invented by adding magnesium to the widely used grey iron. In 1940, thermal sand reclamation was applied for molding and core sands. In 1952, the "D-process" was developed for making shell molds with fine, pre-coated sand. In 1953, the hotbox core sand process in which the cores are thermally cured was invented. In 1954, a new core binder—water glass (sodium silicate), hardened with CO2 from ambient air, came out In the 2010s, additive manufacturing began to be applied to sand mold preparation in commercial production; instead of the sand mold being formed via packing sand around a pattern, it is 3D-printed.	Technology	Metallurgy	null
67656	https://en.wikipedia.org/wiki/Prunus	Prunus	Prunus is a genus of flowering trees and shrubs from the family Rosaceae, which includes plums, cherries, peaches, nectarines, apricots and almonds (collectively stonefruit). The genus has a cosmopolitan distribution, being native to the temperate regions of North America, the neotropics of South America, and temperate and tropical regions of Eurasia and Africa, There are about 340 accepted species . Many members of the genus are widely cultivated for their sweet, fleshy fruit and for decorative purposes of their flowers. Prunus fruit are drupes, or stone fruits. The fleshy mesocarp surrounding the endocarp is edible while the endocarp itself forms a hard, inedible shell called the pyrena ("stone" or "pit"). This shell encloses the seed (or "kernel"), which is edible in some species (such as sweet almonds), but poisonous in many others (such as apricot kernels). Besides being eaten off the hand, most Prunus fruit are also commonly used in processing, such as jam production, canning, drying, and the seeds for roasting. Description Members of the genus are either deciduous or evergreen. A few species have spiny stems. The leaves are simple, alternate, usually lanceolate, unlobed, and often with nectaries on the leaf stalk along with stipules. The flowers are usually white to pink, sometimes red, with five petals and five sepals. Numerous stamens are present. Flowers are borne singly, or in umbels of two to six or sometimes more on racemes. The fruit is a fleshy drupe (a "prune") with a single relatively large, hard-coated seed (a "stone"). Taxonomy Within the rose family Rosaceae, it was traditionally placed as a subfamily, the Amygdaloideae (incorrectly "Prunoideae"), but was sometimes placed in its own family, the Prunaceae (or Amygdalaceae). More recently, Prunus is thought to have evolved from within a much larger clade now called subfamily Amygdaloideae (incorrectly "Spiraeoideae"). Classification Evolutionary history The oldest fossils confirmed to belong to Prunus date to the Eocene, and are found across the Northern Hemisphere. Older potential Late Cretaceous records are unconfirmed. The earliest known fossil Prunus specimens are wood, drupe, seed, and a leaf from the middle Eocene of the Princeton Chert of British Columbia, Canada. Using the known age as calibration data, a partial phylogeny of some of the Rosaceae from a number of nucleotide sequences was reconstructed. Prunus and its sister clade Maloideae (apple subfamily) has been suggested to have diverged which is within the Lutetian, or older middle Eocene. Stockey and Wehr report: "The Eocene was a time of rapid evolution and diversification in Angiosperm families such as the Rosaceae ...." The oldest fossil species is Prunus cathybrownae from the Klondike Mountain Formation. The Princeton finds are among a large number of angiosperm fossils from the Okanagan Highlands dating to the late early and middle Eocene. Crataegus is found at three locations: the McAbee Fossil Beds, British Columbia; the Klondike Mountain Formation around Republic, Washington, and the Allenby Formation around Princeton, British Columbia, while Prunus is found at those locations plus the Coldwater Beds of Quilchena, British Columbia and Chu Chua Formation around Chu Chua, British Columbia. A review of research on the Eocene Okanagan Highlands reported that the Rosaceae were more diverse at higher altitudes. The Okanagan highlands formations date to as early as 52 mya, but the (approximate) 44.3 mya date might still apply. The authors state that "the McAbee flora records a diverse early middle Eocene angiosperm-dominated forest." Linnean classification In 1737, Carl Linnaeus used four genera to include the species of modern Prunus—Amygdalus, Cerasus, Prunus, and Padus—but simplified it to Amygdalus and Prunus in 1758. Since then, the various genera of Linnaeus and others have become subgenera and sections, as all the species clearly are more closely related. Liberty Hyde Bailey said: "The numerous forms grade into each other so imperceptibly and inextricably that the genus cannot be readily broken up into species." Traditional classification Historical treatments break the genus into several different genera, but this segregation is not currently widely recognised other than at the subgeneric rank. The ITIS recognises just the single genus Prunus, with an open list of species, all of which are given at List of Prunus species. One treatment of the subgenera derives from the work of Alfred Rehder in 1940. Rehder hypothesized five subgenera: Amygdalus, Prunus, Cerasus, Padus, and Laurocerasus. To them C. Ingram added Lithocerasus. The six subgenera are described as follows: Subgenus Amygdalus, almonds and peaches: axillary buds in threes (vegetative bud central, two flower buds to sides); flowers in early spring, sessile or nearly so, not on leafed shoots; fruit with a groove along one side; stone deeply grooved; type species: Prunus dulcis (almond) Subgenus Prunus, plums and apricots: axillary buds solitary; flowers in early spring stalked, not on leafed shoots; fruit with a groove along one side, stone rough; type species: Prunus domestica (plum) Subgenus Cerasus, true cherries: axillary buds single; flowers in early spring in corymbs, long-stalked, not on leafed shoots; fruit not grooved, stone smooth; type species: Prunus cerasus (sour cherry) Subgenus Lithocerasus, bush cherries: axillary buds in threes; flowers in early spring in corymbs, long-stalked, not on leafed shoots; fruit not grooved, stone smooth; type species: Prunus pumila (sand cherry) Subgenus Padus, bird cherries: axillary buds single; flowers in late spring in racemes on leafy shoots, short-stalked; fruit not grooved, stone smooth; type species: Prunus padus (European bird cherry), now known to be polyphyletic Subgenus Laurocerasus, cherry laurels: mostly evergreen (all the other subgenera are deciduous); axillary buds single; flowers in early spring in racemes, not on leafed shoots, short-stalked; fruit not grooved, stone smooth; type species: Prunus laurocerasus (European cherry-laurel) Phylogenetic classification An extensive phylogenetic study based on different chloroplast and nuclear sequences divides Prunus into three subgenera: Subg. Padus: In addition to species of Padus (bird cherries), this subgenus also includes species of Maddenia (false bird cherries), Laurocerasus (cherry laurels) and Pygeum. Subg. Cerasus: This subgenus includes true cherries such as sweet cherry, sour cherry, mahaleb cherry and Japanese flowering cherry. Subg. Prunus: This subgenus includes the following sections: Sect. Prunus: Old World plums Sect. Prunocerasus: New World plums Sect. Armeniaca: apricots Sect. Microcerasus: bush cherries Sect. Amygdalus: almonds Sect. Persica: peaches Sect. Emplectocladus: desert almonds Species The lists below are incomplete, but include most of the better-known species. Afro-Eurasian species P. africana – African cherry P. apetala – clove cherry P. armeniaca – apricot P. avium – sweet cherry or wild cherry P. brigantina – Briançon apricot P. buergeriana – dog cherry P. campanulata – Taiwan cherry P. canescens – gray-leaf cherry P. cerasifera – cherry plum P. cerasoides – wild Himalayan cherry P. cerasus – sour cherry P. ceylanica – Ceylon cherry P. cocomilia – Italian plum P. cornuta – Himalayan bird cherry P. davidiana – David's peach P. darvasica – Darvaz plum P. domestica – common plum P. dulcis – almond P. fruticosa – European dwarf cherry P. glandulosa – Chinese bush cherry P. grayana – Japanese bird cherry P. incana – willow-leaf cherry P. incisa – Fuji cherry P. jacquemontii – Afghan bush cherry P. japonica – Japanese bush cherry P. laurocerasus – cherry laurel P. lusitanica – Portugal laurel P. maackii – Manchurian cherry P. mahaleb – Mahaleb cherry P. mandshurica – Manchurian apricot P. maximowiczii – Korean cherry P. mume – Chinese plum P. nipponica – Japanese alpine cherry P. padus – bird cherry P. persica – peach P. pseudocerasus – Chinese sour cherry P. prostrata – mountain cherry P. salicina – Japanese plum P. sargentii – north Japanese hill cherry P. scoparia – mountain almond P. serrula – Tibetan cherry P. serrulata – Japanese cherry P. sibirica – Siberian apricot P. simonii – apricot plum P. speciosa – Oshima cherry P. spinosa – blackthorn, sloe P. ssiori – Hokkaido bird cherry P. subhirtella – winter-flowering cherry P. tenella – dwarf Russian almond P. tomentosa – Nanking cherry P. triloba – flowering plum P. turneriana – almondbark P. ursina – Bear's plum P. × yedoensis – Yoshino cherry P. zippeliana – big-leaf cherry (Chinese: 大叶桂樱) Species found in the Americas P. alabamensis – Alabama cherry P. alleghaniensis – Allegheny plum P. americana – American plum P. andersonii – desert peach P. angustifolia – Chickasaw plum P. brasiliensis – Brazilian cherry P. buxifolia – chuwacá P. caroliniana – Carolina laurelcherry P. cortapico P. emarginata – bitter cherry P. eremophila – Mojave Desert plum P. fasciculata – wild almond P. fremontii – desert apricot P. geniculata – scrub plum P. gentryi – Gentry cherry P. gracilis – Oklahoma plum P. havardii – Havard's plum P. hortulana – Hortulan plum P. huantensis P. ilicifolia – hollyleaf cherry P. integrifolia P. maritima – beach plum P. mexicana – Mexican plum P. minutiflora – Texas almond P. murrayana – Murray's plum P. myrtifolia – West Indies cherry P. nigra – Canada plum P. occidentalis – western cherry laurel P. pensylvanica – pin cherry P. pleuradenia – Antilles cherry P. pumila – sand cherry P. rigida P. rivularis – creek plum P. serotina – black cherry P. subcordata – Klamath plum P. subcorymbosa P. texana – peachbush P. umbellata – flatwoods plum P. virginiana – chokecherry Etymology The Online Etymology Dictionary presents the customary derivations of plum and prune from Latin prūnum, the plum fruit. The tree is prūnus; and Pliny uses prūnus silvestris to mean the blackthorn. The word is not native Latin, but is a loan from Greek προῦνον (), which is a variant of προῦμνον (), origin unknown. The tree is προύμνη (). Most dictionaries follow Hoffman, Etymologisches Wörterbuch des Griechischen, in making some form of the word a loan from a pre-Greek language of Asia Minor, related to Phrygian. The first use of Prunus as a genus name was by Carl Linnaeus in Hortus Cliffortianus of 1737, which went on to become Species Plantarum. Pests and diseases Various Prunus species are winter hosts of the Damson-hop aphid, Phorodon humuli, which is destructive to hops Humulus lupulus just at the time of their maturity, so plum trees should not be grown in the vicinity of hop fields. Corking is the drying or withering of fruit tissue. In stone fruit, it is often caused by a lack of boron and/or calcium. Gummosis is a nonspecific condition of stone fruits (peach, nectarine, plum, and cherry) in which gum is exuded and deposited on the bark of trees. Gum is produced in response to any type of wound – insect, mechanical injury, or disease. Apiosporina morbosa is a major fungal disease in the Northern Americas, with many urban centres running black knot fungus management programs. This disease is best managed by physical removal of knot-bearing branches to prevent spore spread and immediate disposal of infected tissue. Chemical treatment is not largely effective, as trees can easily be re-infected by neighbouring knots. Laetiporus gilbertsoni (commonly sulfur shelf and chicken of the woods), is a serious cubic brown rot parasite which attacks certain species of decorative red-leaf plum trees in the genus Prunus on the Pacific coast of North America. Cultivation The genus Prunus includes the almond, the nectarine and peach, several species of apricots, cherries, and plums, all of which have cultivars developed for commercial fruit and nut production. The almond is not a true nut; the edible part is the seed. Other species are occasionally cultivated or used for their seed and fruit. A number of species, hybrids, and cultivars are grown as ornamental plants, usually for their profusion of flowers, sometimes for ornamental foliage and shape, and occasionally for their bark. Because of their considerable value as both food and ornamental plants, many Prunus species have been introduced to parts of the world to which they are not native, some becoming naturalised. The Tree of 40 Fruit has 40 varieties grafted on to one rootstock. Species such as blackthorn (Prunus spinosa), are grown for hedging, game cover, and other utilitarian purposes. The wood of some species (notably black cherry) is prized as a furniture and cabinetry timber, especially in North America. Many species produce an aromatic gum from wounds in the trunk; this is sometimes used medicinally. Other minor uses include dye production. Pygeum, a herbal remedy containing extracts from the bark of Prunus africana, is used as to alleviate some of the discomfort caused by inflammation in patients with benign prostatic hyperplasia. Prunus species are food plants for the larvae of many Lepidoptera species (butterflies and moths). Prunus species are included in the Tasmanian Fire Service's list of low flammability plants, indicating that it is suitable for growing within a building protection zone. Ornamental Prunus Ornamentals include the group that may be collectively called "flowering cherries" (including sakura, the Japanese flowering cherries). Toxicity Many species are cyanogenic; that is, they contain compounds called cyanogenic glucosides, notably amygdalin, which, on hydrolysis, yield hydrogen cyanide. Although the fruits of some may be edible by humans and livestock (in addition to the ubiquitous fructivory of birds), seeds, leaves and other parts may be toxic, some highly so. The plants contain no more than trace amounts of hydrogen cyanide, but on decomposition after crushing and exposure to air or on digestion, poisonous amounts may be generated. The trace amounts may give a characteristic taste ("bitter almond") with increasing bitterness in larger quantities, less tolerable to people than to birds, which habitually feed on specific fruits. Benefits to human health People are often encouraged to consume many fruits because they are rich in a variety of nutrients and phytochemicals that are supposedly beneficial to human health. The fruits of Prunus often contain many phytochemicals and antioxidants. These compounds have properties that have been linked to preventing different diseases and disorders. Research suggests that the consumption of these fruits reduces the risk of developing diseases such as cardiovascular diseases, cancer, diabetes, and other age-related declines. Many factors can affect the levels of bioactive compounds in the different fruits of the genus Prunus, including the environment, season, processing methods, orchard operations, and postharvest management. Cherries Cherries contain many different phenolic compounds and anthocyanins, which are indicators of being rich in antioxidants. Recent research has linked the phenolic compounds of the sweet cherry (Prunus avium) with antitumor properties. Reactive oxygen species (ROS) include superoxide radicals, hydrogen peroxide, hydroxyl radicals, and singlet oxygen; they are the byproducts of metabolism. High levels of ROS lead to oxidative stress, which causes damage to lipids, proteins, and nucleic acids. The oxidative damage results in cell death, which ultimately leads to numerous diseases and disorders. Antioxidants act as a defense mechanism against the oxidative stress. They are used to remove the free radicals in a living system that are generated as ROS. Some of those antioxidants include gutathione S-transferase, glutathione peroxidase, superoxide dismutase, and catalase. The antioxidants present in cherry extracts act as inhibitors of the free radicals. However, the DNA and proteins can be damaged when an imbalance occurs in the level of free radicals and the antioxidants. When not enough antioxidants are available to remove the free radicals, many diseases can occur, such as cancers, cardiovascular diseases, Parkinson's disease, etc. Recent studies have shown that using natural antioxidants as a supplement in chemotherapy can decrease the amount of oxidative damage. Some of these natural antioxidants include vitamin C, tocopherol, and epigallocatechin gallate; they can be found in certain cherry extracts. Almonds Similar to cherries, strawberries, and raspberries, almonds are also rich in phenolics. Almonds have a high oxygen radical absorbing capacity (ORAC), which is another indicator of being rich in antioxidants. As stated before, high levels of free radicals are harmful, thus having the capacity to absorb those radicals is greatly beneficial. The bioactive compounds, polyphenols and anthocyanins, found in berries and cherries are also present in almonds. Almonds also contain nonflavonoid and flavonoid compounds, which contribute to their antioxidant properties. Flavonoids are a group of structurally related compounds that are arranged in a specific manner and can be found in all vascular plants on land. They also contribute to the antioxidant properties of almonds. Some of the nonflavonoid compounds present are protocatechuic, vanillic, and p-hydroxybenzoic acids. Flavonoid compounds that can be found in the skin of the almond are flavanols, dihydroflavonols, and flavanones. Plums Of all of the different species of stone fruits, plums are the richest in antioxidants and phenolic compounds. The total antioxidant capacity (TAC) varies within each fruit, but in plums, TAC is much higher in the skin than in the flesh of the fruit. Apricots Apricots are high in carotenoids, which play a key role in light absorption during development. Carotenoids are the pigments that give the pulp and peel of apricots and other Prunus fruits their yellow and orange colors. Moreover, it is an essential precursor for vitamin A, which is especially important for vision and the immune system in humans. Moreover, these fruits are quite rich in phenolic substances, including catechin, epicatechin, p-coumaric acid, caffeic acid, and ferulic acid. Peaches and nectarines Similar to the plum, peaches and nectarines also have higher TAC in the skin than in the flesh. They also contain moderate levels of carotenoids and ascorbic acid. Peaches and nectarines are orange and yellow in color, which can be attributed to the carotenoids present.	Biology and health sciences	Rosales	Plants
67668	https://en.wikipedia.org/wiki/Deimos%20%28moon%29	Deimos (moon)	Deimos (; systematic designation: Mars II) is the smaller and outer of the two natural satellites of Mars, the other being Phobos. Deimos has a mean radius of and takes 30.3 hours to orbit Mars. Deimos is from Mars, much farther than Mars's other moon, Phobos. It is named after Deimos, the Ancient Greek god and personification of dread and terror. Discovery and Etymology Deimos was discovered by Asaph Hall at the United States Naval Observatory in Washington, D.C., on 12 August 1877, at about 07:48 UTC. Hall, who also discovered Phobos shortly afterwards, had been specifically searching for Martian moons at the time. The moon is named after Deimos, a figure representing dread in Greek mythology. The name was suggested by academic Henry Madan, who drew from Book XV of the Iliad, where Ares (Greek counterpart of the Roman god Mars) summons Dread (Deimos) and Fear (Phobos). Origin The origin of Mars' moons is unknown and the hypotheses are controversial. The main hypotheses are that they formed either by capture or by accretion. Because of the postulated similarity to the composition of C- or D-type asteroids, one hypothesis is that the moons may be objects captured into Martian orbit from the asteroid belt, with orbits that have been circularized either by atmospheric drag or tidal forces, as capture requires dissipation of energy. The current Martian atmosphere is too thin to capture a Phobos-sized object by atmospheric braking. Geoffrey Landis has pointed out that the capture could have occurred if the original body was a binary asteroid that separated due to tidal forces. The main alternative hypothesis is that the moons accreted in the present position. Another hypothesis is that Mars was once surrounded by many Phobos- and Deimos-sized bodies, perhaps ejected into orbit around it by a collision with a planetesimal. In 2021, Amirhossein Bagheri (ETH Zurich), Amir Khan (ETH Zurich), Michael Efroimsky (US Naval Observatory) and their colleagues proposed a new hypothesis on the origin of the moons. By analyzing the seismic and orbital data from the Mars InSight Mission and other missions, they proposed that the moons were born from the disruption of a common parent body around 1 to 2.7 billion years ago. The common progenitor of Phobos and Deimos was most probably hit by another object and shattered to form Phobos and Deimos. Physical characteristics Deimos is a gray-colored body. Like most bodies of its size, Deimos is highly non-spherical with triaxial dimensions of , corresponding to a mean diameter of which makes it about 57% the size of Phobos. Deimos is composed of rock rich in carbonaceous material, much like C-type asteroids and carbonaceous chondrite meteorites. It is cratered, but the surface is noticeably smoother than that of Phobos, caused by the partial filling of craters with regolith. The regolith is highly porous and has a radar-estimated density of only . Escape velocity from Deimos is 5.6 m/s. This velocity could theoretically be achieved by a human performing a vertical jump. The apparent magnitude of Deimos is 12.45. Named geological features Only two geological features on Deimos have been given names. The craters Swift and Voltaire are named after writers who speculated on the existence of two Martian moons before Phobos and Deimos were discovered. Orbital characteristics Deimos's orbit is nearly circular and is close to Mars's equatorial plane. Deimos is possibly an asteroid that was perturbed by Jupiter into an orbit that allowed it to be captured by Mars, though this hypothesis is still controversial and disputed. Both Deimos and Phobos have very circular orbits which lie almost exactly in Mars's equatorial plane, and hence a capture origin requires a mechanism for circularizing the initially highly eccentric orbit, and adjusting its inclination into the equatorial plane, most likely by a combination of atmospheric drag and tidal forces; it is not clear that sufficient time was available for this to have occurred for Deimos. As seen from Mars, Deimos would have an angular diameter of no more than 2.5 minutes (sixty minutes make one degree), one twelfth of the width of the Moon as seen from Earth, and would therefore appear almost star-like to the naked eye. At its brightest ("full moon") it would be about as bright as Venus is from Earth; at the first- or third-quarter phase it would be about as bright as Vega. With a small telescope, a Martian observer could see Deimos's phases, which take 1.2648 days (Deimos's synodic period) to run their course. Unlike Phobos, which orbits so fast that it rises in the west and sets in the east, Deimos rises in the east and sets in the west, slower than Mars's rotation speed. The Sun-synodic orbital period of Deimos of about 30.4 hours exceeds the Martian solar day ("sol") of about 24.7 hours by such a small amount that 2.48 days (2.41 sols) elapse between its rising and setting for an equatorial observer. From Deimos-rise to Deimos-rise (or setting to setting), 5.466 days (5.320 sols) elapse. Because Deimos's orbit is relatively close to Mars and has only a very small inclination to Mars's equator, it cannot be seen from Martian latitudes greater than 82.7°. Deimos's orbit is slowly getting larger, because it is far enough away from Mars and because of tidal acceleration. It is expected to eventually escape Mars's gravity. Solar transits Deimos regularly passes in front of the Sun as seen from Mars. It is too small to cause a total eclipse, appearing only as a small black dot moving across the Sun. Its angular diameter is only about 2.5 times the angular diameter of Venus during a transit of Venus from Earth. On 4 March 2004 a transit of Deimos was photographed by Mars rover Opportunity, and on 13 March 2004 a transit was photographed by Mars rover Spirit. Exploration Overall, its exploration history is similar to those of Mars and of Phobos. Deimos has been photographed close-up by several spacecraft whose primary mission has been to photograph Mars, including in March 2023 during a rare close encounter by the Emirates Mars Mission. No landings on Deimos have been made. In 1997 and 1998, the proposed Aladdin mission was selected as a finalist in the NASA Discovery Program. The plan was to visit both Phobos and Deimos, and launch projectiles at the satellites. The probe would collect the ejecta as it performed a slow flyby (~1 km/s). These samples would be returned to Earth for study three years later. The principal investigator was Carle M. Pieters of Brown University. The total mission cost, including launch vehicle and operations was $247.7 million. Ultimately, the mission chosen to fly was MESSENGER, a probe to the planet Mercury. In 2008, NASA Glenn Research Center began studying a Phobos and Deimos sample-return mission that would use solar electric propulsion. The study gave rise to the "Hall" mission concept, a New Frontiers-class mission currently under further study. Also, the sample-return mission called Gulliver has been conceptualized and dedicated to Deimos, in which 1 kilogram (2.2 pounds) of material from Deimos would be returned to Earth. Another concept of sample-return mission from Phobos and Deimos is OSIRIS-REx 2, which would use heritage from the first OSIRIS-REx. In March 2014, a Discovery class mission was proposed to place an orbiter in Mars orbit by 2021 and study Phobos and Deimos. It is called Phobos And Deimos & Mars Environment (PADME). Human exploration of Deimos could serve as a catalyst for the human exploration of Mars. Recently, it was proposed that the sands of Deimos or Phobos could serve as a valuable material for aerobraking in the colonization of Mars. See Phobos for more detail. ISRO's Mars Orbiter Mission captured the first pictures of the far side on Deimos. In April 2023, astronomers released close-up global images, for the first time, of Deimos that were taken by the Mars Hope orbiter. Observations reported by this mission contravene the captured asteroid hypothesis and indicate basaltic planetary origin of Deimos.	Physical sciences	Solar System	Astronomy
67679	https://en.wikipedia.org/wiki/Bullet	Bullet	A bullet is a kinetic projectile, a component of firearm ammunition that is shot from a gun barrel. They are made of a variety of materials, such as copper, lead, steel, polymer, rubber and even wax; and are made in various shapes and constructions (depending on the intended applications), including specialized functions such as hunting, target shooting, training, and combat. Bullets are often tapered, making them more aerodynamic. Bullet size is expressed by weight and diameter (referred to as "caliber") in both imperial and metric measurement systems. Bullets do not normally contain explosives but strike or damage the intended target by transferring kinetic energy upon impact and penetration. Description The term bullet is from Early French, originating as the diminutive of the word boulle (boullet), which means "small ball". Bullets are available singly (as in muzzle-loading and cap and ball firearms) but are more often packaged with propellant as a cartridge ("round" of ammunition) consisting of the bullet (i.e., the projectile), the case (which holds everything together), the propellant (which provides the majority of the energy to launch the projectile), and the primer (which ignites the propellant). Cartridges, in turn, may be held in a magazine, a clip, or a belt (for rapid-fire automatic firearms). Although the word bullet is often used in colloquial language to refer to a cartridge round, a bullet is not a cartridge but rather a component of one. This use of the term bullet (when intending to describe a cartridge) often leads to confusion when a cartridge and all its components are specifically being referenced. The sound of gunfire (i.e. the "muzzle report") is often accompanied with a loud bullwhip-like crack as the supersonic bullet pierces through the air, creating a sonic boom. Bullet speeds at various stages of flight depend on intrinsic factors such as sectional density, aerodynamic profile and ballistic coefficient, as well as extrinsic factors such as barometric pressure, humidity, air temperature and wind speed. Subsonic cartridges fire bullets slower than the speed of sound, so there are no sonic booms. This means that a subsonic cartridge, such as .45 ACP, can be substantially quieter than a supersonic cartridge, such as the .223 Remington, even without the use of a suppressor. Bullets shot by firearms can be used for target practice or to injure or kill animals or people. Death can be by blood loss or damage to vital organs, or even asphyxiation if blood enters the lungs. Bullets are not the only projectiles shot from firearm-like equipment: BBs are shot from BB guns, airsoft pellets are shot by airsoft guns, paintballs are shot by paintball markers, and small rocks can be hurtled from slingshots. There are also flare guns, potato guns (and spud guns), tasers, bean bag rounds, grenade launchers, flash bangs, tear gas, RPGs, and missile launchers. Speed Bullets used in many cartridges are fired at muzzle velocities faster than the speed of sound—about in dry air at —and thus can travel substantial distances to their targets before any nearby observers hear the sound of the shots. Rifle bullets, such as that of a Remington 223 firing lightweight varmint projectiles from a 24 inch barrel, leave the muzzle at speeds of up to . A bullet from a 9 mm Luger handgun, reaches speeds of only . Similarly, an AK-47, has a muzzle velocity of about . History The first true gun evolved in China from the fire lance (a bamboo tube that fired porcelain shrapnel) with the invention of the metal hand cannon sometime around 1288, which the Yuan dynasty used to win a decisive victory against Mongolian rebels. The artillery cannon appeared in 1326 and the European hand cannon in 1364. Early projectiles were made of stone. Eventually it was discovered that stone would not penetrate stone fortifications, which led to the use of denser materials as projectiles. Hand cannon projectiles developed in a similar manner. The first recorded instance of a metal ball from a hand cannon penetrating armor was in 1425. Shot retrieved from the wreck of the Mary Rose (sunk in 1545, raised in 1982) are of different sizes, and some are stone while others are cast iron. The development of the hand culverin and matchlock arquebus brought about the use of cast lead balls as projectiles. The original round musket ball was smaller than the bore of the barrel. At first it was loaded into the barrel just resting upon the powder. Later, some sort of material was used as a wadding between the ball and the powder as well as over the ball to keep it in place, it held the bullet firmly in the barrel and against the powder. (Bullets not firmly set on the powder risked exploding the barrel, with the condition known as a "short start".) The loading of muskets was therefore easy with the old smooth-bore Brown Bess and similar military muskets. The original muzzle-loading rifle, however, was loaded with a piece of leather or cloth wrapped around the ball, to allow the ball to engage the grooves in the barrel. Loading was a bit more difficult, particularly when the bore of the barrel was fouled from previous firings. For this reason, and because rifles were not often fitted for bayonets, early rifles were rarely used for military purposes, compared to muskets. There was a distinct change in the shape and function of the bullet during the first half of the 19th century, although experiments with various types of elongated projectiles had been made in Britain, America and France from the first half of the 18th century onwards. In 1816, Capt. George Reichenbach of the Bavarian army invented a rifled-wall musket using cylindro-conical ammunition. In 1826, Henri-Gustave Delvigne, a French infantry officer, invented a breech with abrupt shoulders on which a spherical bullet was rammed down until it caught the rifling grooves. Delvigne's method, however, deformed the bullet and was inaccurate. In 1855, a detachment of 1st U.S. Dragoons, while on patrol, traded lead for gold bullets with Pima Indians along the California–Arizona border. Square bullets have origins that almost pre-date civilization and were used in slings. They were typically made out of copper or lead. The most notable use of square bullet designs was by James Puckle and Kyle Tunis who patented them, where they were briefly used in one version of the Puckle gun. The early use of these in the black-powder era was soon discontinued because of the irregular and unpredictable flight patterns. Pointed bullets Delvigne continued to develop bullet design and by 1830 had started to develop cylindro-conical bullets. His bullet designs were improved by Francois Tamisier with the addition of "ball grooves" which are known as "cannelures", which moved the resistance of air behind the center of gravity of the bullet. Tamisier also developed progressive rifling: the rifle grooves were deeper toward the breech, becoming shallower as they progressed toward the muzzle. This causes the bullet to be progressively molded into the grooves which increases range and accuracy. Among the first pointed or "conical" bullets were those designed by Captain John Norton of the British Army in 1832. Norton's bullet had a hollow base made of lotus pith that on firing expanded under pressure to engage with a barrel's rifling. The British Board of Ordnance rejected it because spherical bullets had been in use for the previous 300 years. Renowned English gunsmith William Greener invented the Greener bullet in 1836. Greener fitted the hollow base of an oval bullet with a wooden plug that more reliably forced the base of the bullet to expand and catch the rifling. Tests proved that Greener's bullet was effective, but the military rejected it because, being two parts, they judged it as too complicated to produce. The carabine à tige, developed by Louis-Étienne de Thouvenin in 1844, was an improvement of Delvigne's design. The rifle barrel has a forcing plug in the breech of the barrel to mold the bullet into the rifling with the use of a special ramrod. While successful in increasing accuracy, it was difficult to clean. The soft lead Minié ball was first introduced in 1847 by Claude-Étienne Minié, a captain in the French Army. It was another improvement of the work done by Delvigne. The bullet was conical in shape with a hollow cavity in the rear, which was fitted with a small iron cap instead of a wooden plug. When fired, the iron cap forced itself into the hollow cavity at the rear of the bullet, thus expanding the sides of the bullet to grip and engage the rifling. In 1851, the British adopted the Minié ball for their 702-inch Pattern 1851 Minié rifle. In 1855, James Burton, a machinist at the U.S. Armory at Harper's Ferry, West Virginia, improved the Minié ball further by eliminating the metal cup in the bottom of the bullet. The Minié ball first saw widespread use in the Crimean War (1853–1856). Roughly 90% of the battlefield casualties in the American Civil War (1861–1865) were caused by Minié balls fired from rifled muskets. A similar bullet called the Nessler ball was also developed for smoothbore muskets. Between 1854 and 1857, Sir Joseph Whitworth conducted a long series of rifle experiments and proved, among other points, the advantages of a smaller bore and, in particular, of an elongated bullet. The Whitworth bullet was made to fit the grooves of the rifle mechanically. The Whitworth rifle was never adopted by the government, although it was used extensively for match purposes and target practice between 1857 and 1866. In 1861, W. B. Chace approached President Abraham Lincoln with an improved ball design for muskets. In firing over the Potomac River, where the Chace ball and the round ball were alternated, Lincoln observed that the Chace design carried a third or more farther fired at the same elevation. Although Lincoln recommended testing, it never took place. Around 1862, W. E. Metford carried out an exhaustive series of experiments on bullets and rifling, and he invented the important system of light rifling with increasing spiral and a hardened bullet. The combined result was that, in December 1888, the Lee–Metford small-bore (.303", 7.70 mm) rifle, Mark I, was adopted for the British army. The Lee–Metford was the predecessor of the Lee–Enfield. Modern bullets The next important change in the history of the rifle bullet occurred in 1882, when Lieutenant Colonel Eduard Rubin, director of the Swiss Army Laboratory at Thun, invented the copper-jacketed bullet — an elongated bullet with a lead core in a copper jacket. It was also small bore (7.5 and 8 mm) and it is the precursor of the 8 mm Lebel bullet adopted for the smokeless powder ammunition of the Lebel Model 1886 rifle. The surface of lead bullets fired at high velocity may melt from the hot gases behind and friction within the bore. Because copper has a higher melting point, and greater specific heat capacity, and higher hardness, copper-jacketed bullets allow greater muzzle velocities. European advances in aerodynamics led to the pointed spitzer bullet. By the beginning of the 20th century, most world armies had begun the transition to spitzer bullets. These bullets flew for greater distances more accurately and transferred more kinetic energy. Spitzer bullets combined with machine guns greatly increased lethality on the battlefield. Spitzer bullets were streamlined at the base with the boat tail. In the trajectory of a bullet, as air passes over a bullet at high speed, a vacuum is created at the end of the bullet, slowing the projectile. The streamlined boat tail design reduces this form drag by allowing the air to flow along the surface of the tapering end. The resulting aerodynamic advantage is currently seen as the optimum shape for rifle technology. The first combination spitzer and boat-tail bullet, named balle D by its inventor Captain Georges Desaleux, was introduced as standard military ammunition in 1901, for the French Lebel Model 1886 rifle. A ballistic tip bullet is a hollow-point rifle bullet that has a plastic tip on the end of the bullet. This improves external ballistics by streamlining the bullet, allowing it to cut through the air more easily, and improves terminal ballistics by allowing the bullet to act as a jacketed hollow point. As a side effect, it also feeds better in weapons that have trouble feeding rounds that are not full metal jacket rounds. Design Bullet designs have to solve two primary problems. In the barrel, they must first form a seal with the gun's bore. If a strong seal is not achieved, gas from the propellant charge leaks past the bullet, thus reducing efficiency and possibly accuracy. The bullet must also engage the rifling without damaging or excessively fouling the gun's bore and without distorting the bullet, which will also reduce accuracy. Bullets must have a surface that forms this seal without excessive friction. These interactions between bullet and bore are termed internal ballistics. Bullets must be produced to a high standard, as surface imperfections can affect firing accuracy. The physics affecting the bullet once it leaves the barrel is termed external ballistics. The primary factors affecting the aerodynamics of a bullet in flight are the bullet's shape and the rotation imparted by the rifling of the gun barrel. Rotational forces stabilize the bullet gyroscopically as well as aerodynamically. Any asymmetry in the bullet is largely canceled as it spins. However, a spin rate greater than the optimum value adds more trouble than good, by magnifying the smaller asymmetries or sometimes resulting in the bullet breaking apart in flight. With smooth-bore firearms, a spherical shape is optimal because no matter how the bullet is oriented, its aerodynamics are similar. These unstable bullets tumble erratically and provide only moderate accuracy; however, the aerodynamic shape changed little for centuries. Generally, bullet shapes are a compromise between aerodynamics, interior ballistic necessities, and terminal ballistics requirements. Terminal ballistics and stopping power are aspects of bullet design that affect what happens when a bullet impacts with an object. The outcome of the impact is determined by the composition and density of the target material, the angle of incidence, and the velocity and physical characteristics of the bullet. Bullets are generally designed to penetrate, deform, or break apart. For a given material and bullet, the strike velocity is the primary factor that determines which outcome is achieved. Bullet shapes are many and varied. With a mold, bullets can be made at home for reloading ammunition, where local laws allow. Hand-casting, however, is only time- and cost-effective for solid lead bullets. Cast and jacketed bullets are also commercially available from numerous manufacturers for handloading and are most often more convenient than casting bullets from bulk or scrap lead. Propulsion Propulsion of the ball can happen via several methods: by using only gunpowder (as in flintlock, wheellock, or matchlock weapons) by using a percussion cap and gunpowder (as in percussion weapons) by using a cartridge Materials Bullets for black powder, or muzzle-loading firearms, were classically molded from pure lead. This worked well for low-speed bullets, fired at velocities of less than 450 m/s (1,475 ft/s). For slightly higher-speed bullets fired in modern firearms, a harder alloy of lead and tin or typesetter's lead (used to mold linotype) works very well. For even higher-speed bullet use, jacketed lead bullets are used. The common element in all of these, lead, is widely used because it is very dense, thereby providing a high amount of mass—and thus, kinetic energy—for a given volume. Lead is also cheap, easy to obtain, easy to work, and melts at a low temperature, which results in comparatively easy fabrication of bullets. Lead: simple cast, extruded, swaged, or otherwise fabricated lead slugs are the simplest form of bullets. At speeds of greater than 300 m/s (1,000 ft/s) (common in most handguns), lead is deposited in rifled bores at an ever-increasing rate. Alloying the lead with a small percentage of tin and/or antimony serves to reduce this effect but grows less effective as velocities are increased. A cup made of harder metal, such as copper, placed at the base of the bullet and called a gas check, is often used to decrease lead deposits by protecting the rear of the bullet against melting when fired at higher pressures, but this does not solve the problem at higher velocities. A modern solution is to powder coat the lead projectile, encasing it in a protective skin, allowing higher velocities to be achieved without lead deposits. Jacketed lead: bullets intended for even higher-velocity applications generally have a lead core that is jacketed or plated with gilding metal, cupronickel, copper alloys, or steel; a thin layer of harder metal protects the softer lead core when the bullet is passing through the barrel and during flight, which allows delivering the bullet intact to the target. There, the heavy lead core delivers its kinetic energy to the target. Full metal jacket or "ball" bullets (cartridges with ball bullets, which despite the name are not spherical, are called ball ammunition) are completely encased in the harder metal jacket, except for the base. Some bullet jackets do not extend to the front of the bullet, to aid expansion and increase lethality; these are called soft point (if the exposed lead tip is solid) or hollow point bullets (if a cavity or hole is present). Steel bullets are often plated with copper or other metals for corrosion resistance during long periods of storage. Synthetic jacket materials such as nylon and Teflon have been used, with limited success, especially in rifles; however, hollow point bullets with plastic aerodynamic tips have been very successful at both improving accuracy and enhancing expansion. Newer plastic coatings for handgun bullets, such as Teflon-coated bullets, are making their way into the market. Solid or monolithic solid: mono-metal bullets intended for deep penetration in big game animals and slender shaped very-low-drag projectiles for long range shooting are produced out of metals like oxygen-free copper and alloys like cupronickel, tellurium copper and brass (e.g., highly machinable UNS C36000 free-cutting brass). Often these projectiles are turned on precision CNC lathes. In the case of solids, and the ruggedness of the game animals on which they are used, e.g., the African buffalo or elephant, expansion is almost entirely relinquished for the necessary penetration. In shotgunning, "slug" loads are often solid large single lead projectiles, sometimes with a hollow point, used for deer or wild pig hunting in jurisdictions that do not allow hunting with rifles (because a missed slug shot will travel considerably less far than a rifle bullet). Fluted: in appearance, these are solid bullets with scalloped sides (missing material). The theory is that the flutes produce hydraulic jetting when passing through tissue, creating a wound channel larger than that made by conventional expanding ammunition such as hollow point bullets. Hard cast: a hard lead alloy intended to reduce fouling of rifling grooves (especially of the polygonal rifling used in some popular pistols). Benefits include simpler manufacture than jacketed bullets and good performance against hard targets; limitations are an inability to mushroom and subsequent over-penetration of soft targets. Blank: wax, paper, plastic, and other materials are used to simulate live gunfire and are intended only to hold the powder in a blank cartridge and to produce noise, flame and smoke. The "bullet" may be captured in a purpose-designed device or it may be allowed to expend what little energy it has in the air. Some blank cartridges are crimped or closed at the end and do not contain any bullet; some are fully loaded cartridges (without bullets) designed to propel rifle grenades. The force of the expanding gas from blank cartridges can be lethal at short range; fatal accidents have occurred with blank cartridges (e.g., the death of actor Jon-Erik Hexum). Practice: made from lightweight materials like rubber, wax, wood, plastic, or lightweight metal, practice bullets are intended for short-range target work only. Because of their weight and low velocity, they have limited range. Polymer: these are metal-polymer composites, generally lighter and having higher velocities than pure metal bullets of the same dimensions. They permit unusual designs that are difficult with conventional casting or lathing. Less lethal, or less than lethal: Rubber bullets, plastic bullets, and beanbags are designed to be non-lethal, e.g., for use in riot control. They are generally low velocity and are fired from shotguns, grenade launchers, paint ball guns, or specially designed firearms and air gun devices. Incendiary: these bullets are made with explosive or flammable mixtures in the tips that are designed to ignite on contact with a target. The intent is to ignite fuel or munitions in the target area, thereby adding to the destructive power of the bullet. Exploding: similar to the incendiary bullet, this type of projectile is designed to explode upon hitting a hard surface, preferably the bone of the intended target. Not to be mistaken for cannon shells or grenades with fuse devices, these bullets have only cavities filled with a small amount of high explosive depending on the velocity and deformation upon impact to detonate. Exploding bullets have been used in various heavy machine guns and in anti-materiel rifles. Tracer: these have hollow backs, filled with a flare material. Usually this is a mixture of magnesium, a perchlorate, and strontium salts to yield a bright red color, although other materials providing other colors have also sometimes been used. Tracer material burns out after a certain amount of time. This allows the shooter to visually trace the flight path of the projectile and thus make necessary ballistic corrections, without having to confirm projectile impacts and without even using the sights of the weapon. This type of round is also used by all branches of the United States military in combat environments as a signaling device to friendly forces. Normally it is loaded at a four to one ratio with ball ammunition. Armor-piercing: jacketed designs where the core material is a very hard, high-density metal such as tungsten, tungsten carbide, depleted uranium, or steel. A pointed tip is often used, but a flat tip on the penetrator portion is generally more effective. Nontoxic shot: steel, bismuth, tungsten, and other alloys prevent release of toxic lead into the environment. Regulations in several countries mandate the use of nontoxic projectiles especially when waterfowl hunting. It has been found that birds swallow small lead shot for their gizzards to grind food (as they would swallow pebbles of similar size), and the effects of lead poisoning by grinding of lead pellets against food means lead poisoning effects are magnified. Such concerns apply primarily to shotguns firing pellets (shot) and not bullets, but there is evidence suggesting that consumption of spent rifle and pistol ammunition is also hazardous to wildlife. Reduction of hazardous substances legislation has also been applied to bullets on occasion to reduce the impact of lead on the environment at shooting ranges. Blended-metal: bullets made using cores from powdered metals other than lead with binder or sometimes sintered. Frangible: designed to disintegrate into tiny particles upon impact to minimize their penetration for reasons of range safety, to limit environmental impact, or to limit the shoot-through danger behind the intended target. An example is the Glaser Safety Slug, usually a pistol caliber bullet made from an amalgam of lead shot and a hard (and thus frangible) plastic binder designed to penetrate a human target and release its component shot pellets without exiting the target. Multiple projectile: bullets that are made of separate slugs that fit together inside the cartridge and act as a single projectile inside the barrel as they are fired. The projectiles part in flight but are held in formation by tethers that keep the individual parts of the "bullet" from flying too far away from each other. The intention of such ammo is to increase hit chance by giving a shot-like spread to rifled slug firing guns, while maintaining a consistency in shot groupings. Multiple impact bullets may be less stable in flight than conventional solid bullets because of the added drag from the tether line holding the pieces in formation, and each projectile affects the flight of all the others. This may limit the benefit provided by the spread of each bullet at longer ranges. Expanding bullets are designed to increase in diameter upon impact with a target, maximizing the transfer of energy and creating a larger wound channel. These bullets are often made with a lead core and a copper jacket, though variations like MRX bullets have tungsten in its core. The polymer tip in expanding bullets is designed to enhance aerodynamics for shooting at flat long-range trajectories. Treaties and prohibitions Poisonous bullets were a subject to an international agreement as early as the Strasbourg Agreement (1675). The Saint Petersburg Declaration of 1868 prohibited the use of explosive projectiles weighing less than 400 grams. The Hague Conventions prohibits certain kinds of ammunition for use in war. These include poisoned and expanding bullets. Protocol III of the 1983 Convention on Certain Conventional Weapons, an annexed protocol to the Geneva Conventions, prohibits the use of incendiary ammunitions against civilians. Types of bullets Some types of bullets include: Armor piercing (sometimes with a depleted uranium or other heavy metal core) Armor-piercing fin stabilized discarding sabot round Cast Expanding (hollow point, soft point) Frangible Full metal jacket (also known as "ball" ammunition) Hollow-base Hollow-point Hydra-Shok Nosler partition Plastic-tipped Sabot Saboted light armor penetrator Spitzer Semiwadcutter Total metal jacket Very low drag Wadcutter Wax	Technology	Ammunition	null
67717	https://en.wikipedia.org/wiki/American%20robin	American robin	The American robin (Turdus migratorius) is a migratory bird of the true thrush genus and Turdidae, the wider thrush family. It is named after the European robin because of its reddish-orange breast, though the two species are not closely related, with the European robin belonging to the Old World flycatcher family. The American robin is widely distributed throughout North America, wintering from southern Canada to central Mexico and along the Pacific coast. According to the Partners in Flight database (2019), the American robin is the most abundant landbird in North America (with 370million individuals), ahead of red-winged blackbirds, introduced European starlings, mourning doves and house finches. It has seven subspecies. The species is active mostly during the day and assembles in large flocks at night. Its diet consists of invertebrates (such as beetle grubs, earthworms, and caterpillars), fruits, and berries. It is one of the earliest bird species to lay its eggs, beginning to breed shortly after returning to its summer range from its winter range. The robin's nest consists of long coarse grass, twigs, paper, and feathers, and is smeared with mud and often cushioned with grass or other soft materials. It is among the earliest birds to sing at dawn, and its song consists of several discrete units that are repeated. The adult's main predator is the domestic cat; other predators include hawks and snakes. When feeding in flocks, it can be vigilant, watching other birds for reactions to predators. Brown-headed cowbirds (Molothrus ater) lay their eggs in robin nests (see brood parasite), but the robins usually reject the egg. Taxonomy This species was first described in 1766 by Carl Linnaeus in the twelfth edition of his Systema Naturae as Turdus migratorius. The binomial name derives from two Latin words: , "thrush", and from "to migrate". The term for this species has been recorded since at least 1703. A 2020 genetic study has shown that the American robin is closest to the rufous-collared thrush (T. rufitorques) of Central America, confirming a 2007 study which also placed this as its closest relative. Though having distinct plumage, the two species are similar in vocalization and behavior. Beyond this, it lies in a group of other Central American thrushes, suggesting a recent spread northwards into North America; the 2007 study suggested rufous-backed thrush (T. rufopalliatus; not included in the 2020 study) as the next closest relative, with both studies giving the next-closest relatives beyond this trio as the species pair of black thrush (T. infuscatus) and sooty thrush (T. nigrescens), also of Central America. These results contrast markedly with two older studies of the mitochondrial cytochrome b gene which had suggested, though only with weak support, that the American robin might be more closely related to the Kurrichane thrush (T. libonyanus) and the olive thrush (T. olivaceus), both African species, rather than other American thrushes. Subspecies Seven subspecies are accepted. These, except for the isolated T. m. confinis, intergrade with each other and are only weakly defined. T. m. nigrideus breeds from coastal northern Quebec to Labrador and Newfoundland and winters from southern Newfoundland south through most of the eastern U.S. states to southern Louisiana, southern Mississippi, and northern Georgia. It is uniformly darker or blackish on the head, with a dark gray back. The underparts are slightly richer red than those of the nominate subspecies. T. m. migratorius, the nominate subspecies, breeds in the U.S. and Canada, other than down the West Coast, to the edge of the tundra from Alaska and northern Canada east to New England and then south to Maryland, northwestern Virginia, and North Carolina. It winters in southern coastal Alaska, southern Canada, most of the U.S., Bermuda, the Bahamas and eastern Mexico. T. m. achrusterus breeds from southern Oklahoma east to Maryland and western Virginia and south to northern Florida and the Gulf Coast states. It winters through much of the southern part of the breeding range. It is marginally smaller than the eastern subspecies. The black feathers of the forehead and crown have pale gray tips. The underparts are paler than those of the eastern subspecies. T. m. caurinus breeds in southeastern Alaska through coastal British Columbia to Washington and northwestern Oregon. It winters from southwestern British Columbia south to central and southern California and east to northern Idaho. It is slightly smaller than the eastern subspecies and very dark-headed. The white on the tips of the outer two tail feathers is restricted. T. m. propinquus breeds from southeastern British Columbia, southern Alberta, and southwestern Saskatchewan south to southern California and northern Baja California. It winters throughout much of the southern breeding range and south to Baja California. It is the same size as, or slightly larger than, the eastern subspecies, but paler and tinged more heavily brownish-gray. It has very little white on the tip of the outermost tail feathers. Some birds, probably females, lack almost any red below. Males are usually darker and may show pale or whitish sides to the head. T. m. phillipsi is resident in Mexico south to central Oaxaca. It is slightly smaller than the western subspecies, but has a larger bill; the male's underparts are less brick-red than the eastern subspecies and have a rustier tone. T. m. confinis breeds above in the Sierra de la Laguna mountains of southern Baja California. This isolated non-migratory subspecies is particularly distinctive. It is relatively small, and the palest subspecies, with uniform pale gray-brown on the head, face, and upperparts, and pale buffy orange underparts. It usually lacks any white spots to the tips of the outer tail feathers, which have white edges. It has sometimes been classed as a separate species, but both the American Ornithologists' Union and the IOC World Bird List regard it as only a subspecies, albeit in a different group from the other six subspecies. Description The eastern subspecies (T. m. migratorius) is long with a wingspan ranging from , with similar size ranges across all subspecies. The species averages about in weight, with males ranging from and females ranging from . Among standard measurements, the wing chord is , the culmen is and the tarsus is . The head varies from jet black to gray, with white eye arcs and white supercilia. The throat is white with black streaks, and the belly and undertail coverts are white. The adult has a brown back and a reddish-orange breast, varying from a rich red maroon to peachy orange. The bill is mainly yellow with a variably dark tip, the dusky area becoming more extensive in winter, and the legs and feet are brown. The sexes are similar, but females tend to be duller in color than males, with a brown tint to the head, brown upperparts, and less-bright underparts. However, some birds cannot be accurately sexed on the sole basis of plumage. Juveniles are paler in color than adult males and have dark spots on their breasts and whitish wing coverts. First-year birds are not easily distinguishable from adults, but they tend to be duller, with first-year males resembling adult females, and a small percentage retain a few juvenile wing coverts or other feathers. Distribution and habitat The species breeds throughout most of North America, from Alaska and Canada southward to northern Florida and Mexico. While robins occasionally overwinter in the northern part of the United States and southern Canada, most migrate to winter south of Canada from Florida and the Gulf Coast to central Mexico, as well as along the Pacific Coast. Most depart south by the end of August and begin to return north in February and March (exact dates vary with latitude and climate). The distance by which they migrate varies significantly depending on their initial habitat; a study found that individual robins tagged in Alaska are known to travel as much as 3.5 times further across seasons than robins tagged in Massachusetts. The species is a rare vagrant to western Europe, where the majority of records have been in Great Britain, where 29 had been recorded up to the end of 2022. The species has occurred as a vagrant to Greenland, Jamaica, Hispaniola, Puerto Rico and Belize. Vagrants to Europe, where identified to subspecies, are the eastern subspecies (T. m. migratorius), but the Greenland birds included the Newfoundland subspecies (T. m. nigrideus), and some of the southern overshots may have been the southern subspecies (T. m. achrusterus). The breeding habitat is woodland and more open farmland and urban areas. It becomes less common as a breeder in the southernmost part of the Deep South of the United States and there prefers large shade trees on lawns. Its winter habitat is similar but includes more open areas. Diseases The species is a known reservoir (carrier) for West Nile virus spread by Culex mosquitoes. While crows and jays are often the first noticed deaths in an area with West Nile virus, the American robin is suspected to be a key host and holds a larger responsibility for the transmission of the virus to humans. This is because, while crows and jays die quickly from the virus, the American robin survives the virus longer, hence spreading it to more mosquitoes, which then transmit the virus to humans and other species. A successful West Nile virus vaccine has been administered to six 3-5 week old American robins. A DNA vaccine injected intramuscularly resulted in a 400-fold decrease in average viral load that would likely make robins noninfectious and unable to spread disease. An oral bait is the preferred method of distribution of the vaccine as it would be easier and cheaper than intramuscular injection, but more research would be needed as the existing formulation did not work orally. Behavior The American robin is active mostly during the day, and on its winter grounds, it assembles in large flocks at night to roost in trees in secluded swamps or dense vegetation. The flocks break up during the day when the birds feed on fruits and berries in smaller groups. During the summer, males defend a breeding territory and are less social. Diet The diet generally consists of around 40 percent small invertebrates (mainly insects), such as earthworms, beetle grubs, caterpillars, and grasshoppers, and 60 percent wild and cultivated fruits and berries. Their ability to switch to berries allows them to winter much farther north than most other North American thrushes. They will flock to fermented Pyracantha berries, and after eating sufficient quantities will exhibit intoxicated behavior, such as falling over while walking. Robins forage primarily on the ground for soft-bodied invertebrates, and find worms by sight (and sometimes by hearing), pouncing on them and then pulling them up. Nestlings are fed mainly on earthworms and other soft-bodied animal prey. In some areas, robins, particularly of the northwestern subspecies (T. m. caurinus), will feed on beaches, taking insects and small mollusks. American robins are common pests of fruit orchards in North America. Due to their insectivorous and frugivorous diet they have evolved to lose sucrase. Sucrose is unpalatable to them and can be used by humans as a deterrent. The species uses auditory, visual, olfactory and possibly vibrotactile cues to find prey, but vision is the predominant mode of prey detection. It is frequently seen running across lawns picking up earthworms, and its running and stopping behavior is a distinguishing characteristic. In addition to hunting visually, it also has the ability to hunt by hearing. Experiments have shown that it can find earthworms underground by simply using its listening skills. It typically will take several short hops and then cock its head left, right or forward to detect movement of its prey. In urban areas, robins will gather in numbers soon after lawns are mowed or where sprinklers are in use. Threats Juveniles and eggs are preyed upon by squirrels, snakes, and some birds. Adults are primarily taken by Accipiter hawks, cats, and larger snakes such as rat snakes and gopher snakes. Canids such as foxes and dogs take fledglings from the ground. Raccoons often prey upon nests, while small agile carnivores such as American martens, ring-tailed cats and long-tailed weasels hunt adults. The greatest predatory impact is probably from raptorial birds. 28 raptorial bird species hunt American robins. Adult robins are most vulnerable while breeding activities, whereas feeding flocks are vigilant for predators. The American robin rejects cowbird eggs, so brood parasitism by the brown-headed cowbird is rare, and the parasite's chick does not often survive to fledging. In a study of 105 juvenile robins, 77.1% were infected with endoparasites, Syngamus sp. being the most commonly encountered, in 57.1% of the birds. Breeding Breeding begins shortly after the returning to the summer range. The species is one of the first North American birds to lay eggs, and normally has two to three broods per breeding season, which lasts from April to July. The nest is most commonly located above the ground in a dense bush or in a fork between two tree branches, and is built by the female alone. The outer foundation consists of long coarse grass, twigs, paper, and feathers. This is lined with smeared mud and cushioned with fine grass or other soft materials. The American robin builds a new nest for each brood; in northern areas the nest for the first clutch will usually be located in an evergreen tree or shrub, while later broods are raised in deciduous trees. The species is not shy about nesting close to human habitations. A clutch consists of three to five cyan-colored eggs, and is incubated by the female alone. The eggs hatch after 14 days, and the chicks leave the nest a further two weeks later. The altricial chicks are naked and have their eyes closed for the first few days after hatching. The chicks are fed earthworms, insects, and berries. Waste accumulation does not occur in the nest because the adults collect and take it away. Chicks are fed, and then raise tails for elimination of waste, a solid white clump that is collected by a parent prior to flying off. All chicks in the brood leave the nest within two days of each other. Juveniles become capable of sustained flight two weeks after fledging. Chicks become sexually mature at one year of age. Bird banders have found that only 25% of young robins survive their first year. The longest known lifespan of an American robin in the wild is 14 years; the average lifespan is about two years. Vocalization The male, as with many thrushes, has a complex and almost continuous song. It is commonly described as a cheery carol, made up of discrete units, often repeated, and spliced together into a string with brief pauses in between. The song varies regionally, and its style varies by the time of day. The song period is from late February or early March to late July or early August; some birds, particularly in the east, sing occasionally into September or later. They are often among the first songbirds to sing as dawn rises or hours before, and last as evening sets in. It usually sings from a high perch in a tree. The song of the San Lucas subspecies is weaker than that of the eastern subspecies and lacks any clear notes. The male sings when storms approach and again when storms have passed. In addition to its song, the species has a number of calls used for communicating specific information, such as when a ground predator approaches and when a nest or another American robin is being directly threatened. Even during nesting season, when they exhibit mostly competitive and territorial behavior, they may still band together to drive away a predator. Conservation status The species has an extensive range, estimated at , and a large population of about 370 million individuals. The western subspecies (T. m. propinquus) in central California is considered to be expanding its range, as is likely the case elsewhere in the United States. It is threatened by climate change and severe weather, but the population trend appears to be stable, and the species does not approach the vulnerable species thresholds under the population trend criterion (>30% decline over ten years or three generations), and therefore International Union for Conservation of Nature evaluated it as least concern. At one point, the bird was killed for its meat, but it is now protected throughout its range in the United States by the Migratory Bird Treaty Act. In culture The American robin is the state bird of Connecticut, Michigan, and Wisconsin. It was depicted on the 1986 Birds of Canada series Canadian $2 note (this note was subsequently withdrawn.) It has a place in Native American mythology. The story of how the robin got its red breast by fanning the dying flames of a campfire to save a Native American man and a boy is similar to those that surround the European robin. The Tlingit people of northwestern North America held it to be a culture hero created by Raven to please the people with its song. The Peace Bridge robins were a family of American robins that attracted minor publicity in the mid-1930s for their prominent nest on the Canadian side of the Peace Bridge connecting Buffalo, New York, to Fort Erie, Ontario. American popular songs featuring this bird include "When the Red, Red Robin (Comes Bob, Bob, Bobbin' Along)", written by Harry M. Woods. Although the comic book superhero Robin was inspired by an N. C. Wyeth illustration of Robin Hood, a later version had his mother nicknaming him Robin because he was born on the first day of spring. The species is considered a symbol of spring. A well-known example is a poem by Emily Dickinson titled "I Dreaded That First Robin So". Among other 19th-century poems about the first robin of spring is "The First Robin" by William Henry Drummond, which, according to the author's wife, is based on a Quebec superstition that whoever sees the first robin of spring will have good luck. The association has continued down to the present day, as, for example, in one Calvin and Hobbes cartoon from 1990 that had Calvin celebrating his inevitable wealth and fame after seeing the first robin of spring. The harbinger of spring sobriquet is borne out by the fact that American robins tend to follow the isotherm north in spring, but also south in fall. In a study of 209 psychology students at the University of California, Berkeley, Eleanor Rosch found that the robin was, in the students' minds, the most prototypical example of a bird (though the students did not have the opportunity to specify the species of robin). Robin egg blue is a color named after the color of the bird's eggs. Gallery	Biology and health sciences	Passerida	null
67728	https://en.wikipedia.org/wiki/Petrochemical	Petrochemical	Petrochemicals (sometimes abbreviated as petchems) are the chemical products obtained from petroleum by refining. Some chemical compounds made from petroleum are also obtained from other fossil fuels, such as coal or natural gas, or renewable sources such as maize, palm fruit or sugar cane. The two most common petrochemical classes are olefins (including ethylene and propylene) and aromatics (including benzene, toluene and xylene isomers). Oil refineries produce olefins and aromatics by fluid catalytic cracking of petroleum fractions. Chemical plants produce olefins by steam cracking of natural gas liquids like ethane and propane. Aromatics are produced by catalytic reforming of naphtha. Olefins and aromatics are the building-blocks for a wide range of materials such as solvents, detergents, and adhesives. Olefins are the basis for polymers and oligomers used in plastics, resins, fibers, elastomers, lubricants, and gels. Global ethylene production was 190 million tonnes and propylene was 120 million tonnes in 2019. Aromatics production is approximately 70 million tonnes. The largest petrochemical industries are located in the United States and Western Europe; however, major growth in new production capacity is in the Middle East and Asia. There is substantial inter-regional petrochemical trade. Primary petrochemicals are divided into three groups depending on their chemical structure: Olefins includes ethene, propene, butenes and butadiene. Ethylene and propylene are important sources of industrial chemicals and plastics products. Butadiene is used in making synthetic rubber. Aromatics includes benzene, toluene and xylenes, as a whole referred to as BTX and primarily obtained from petroleum refineries by extraction from the reformate produced in catalytic reformers using naphtha obtained from petroleum refineries. Alternatively, BTX can be produced by aromatization of alkanes. Benzene is a raw material for dyes and synthetic detergents, and benzene and toluene for isocyanates MDI and TDI used in making polyurethanes. Manufacturers use xylenes to produce plastics and synthetic fibers. Synthesis gas is a mixture of carbon monoxide and hydrogen used to produce methanol and other chemicals. Steam crackers are not to be confused with steam reforming plants used to produce hydrogen for ammonia production. Ammonia is used to make the fertilizer urea and methanol is used as a solvent and chemical intermediate. Methane, ethane, propane and butanes obtained primarily from natural gas processing plants. Methanol and formaldehyde. In 2007, the amounts of ethylene and propylene produced in steam crackers were about 115 Mt (megatonnes) and 70 Mt, respectively. The output ethylene capacity of large steam crackers ranged up to as much as 1.0 – 1.5 Mt per year. The adjacent diagram schematically depicts the major hydrocarbon sources and processes used in producing petrochemicals. Like commodity chemicals, petrochemicals are made on a very large scale. Petrochemical manufacturing units differ from commodity chemical plants in that they often produce a number of related products. Compare this with specialty chemical and fine chemical manufacture where products are made in discrete batch processes. Petrochemicals are predominantly made in a few manufacturing locations around the world, for example in Jubail and Yanbu Industrial Cities in Saudi Arabia, Texas and Louisiana in the US, in Teesside in the Northeast of England in the United Kingdom, in Tarragona in Catalonia, in Rotterdam in the Netherlands, in Antwerp in Belgium, in Jamnagar, Dahej in Gujarat, India and in Singapore. Not all of the petrochemical or commodity chemical materials produced by the chemical industry are made in one single location but groups of related materials are often made in adjacent manufacturing plants to induce industrial symbiosis as well as material and utility efficiency and other economies of scale. This is known in chemical engineering terminology as integrated manufacturing. Specialty and fine chemical companies are sometimes found in similar manufacturing locations as petrochemicals but, in most cases, they do not need the same level of large-scale infrastructure (e.g., pipelines, storage, ports, and power, etc.) and therefore can be found in multi-sector business parks. The large-scale petrochemical manufacturing locations have clusters of manufacturing units that share utilities and large-scale infrastructures such as power stations, storage tanks, port facilities, road and rail terminals. In the United Kingdom, for example, there are four main locations for such manufacturing: near the River Mersey in North West England, on the Humber on the East coast of Yorkshire, in Grangemouth near the Firth of Forth in Scotland, and in Teesside as part of the Northeast of England Process Industry Cluster (NEPIC). To demonstrate the clustering and integration, some 50% of the United Kingdom's petrochemical and commodity chemicals are produced by the NEPIC industry cluster companies in Teesside. History In 1835, Henri Victor Regnault, a French chemist left vinyl chloride in the sun and found white solid at the bottom of the flask which was polyvinyl chloride. In 1839, Eduard Simon discovered polystyrene by accident by distilling storax. In 1856, William Henry Perkin discovered the first synthetic dye, Mauveine. In 1888, Friedrich Reinitzer, an Austrian plant scientist observed cholesteryl benzoate had two different melting points. In 1909, Leo Hendrik Baekeland invented bakelite made from phenol and formaldehyde. In 1920, Union Carbide built in West Virginia first petrochemical plant in the world. In 1928, synthetic fuels were invented using Fischer-Tropsch process. In 1929, Walter Bock invented synthetic rubber Buna-S which is made up of styrene and butadiene and used to make car tires. In 1933, Otto Röhm polymerized the first acrylic glass methyl methacrylate. In 1935, Michael Perrin invented polyethylene. In 1937, Wallace Hume Carothers invented nylon. In 1938, Otto Bayer invented polyurethane. In 1941, Roy Plunkett invented Teflon. In 1946, he invented Polyester. Polyethylene terephthalate (PET) bottles are made from ethylene and paraxylene. In 1949, Fritz Stastny turned polystyrene into foam. After World War II, polypropylene was discovered in the early 1950s. In 1965, Stephanie Kwolek invented Kevlar. Olefins The following is a partial list of major commercial petrochemicals and their derivatives: ethylene – the simplest olefin; used as a chemical feedstock and ripening stimulant polyethylene – polymerized ethylene; LDPE, HDPE, LLDPE ethanol – via ethylene hydration (chemical reaction adding water) of ethylene ethylene oxide – via ethylene oxidation ethylene glycol – via ethylene oxide hydration engine coolant – ethylene glycol, water and inhibitor mixture polyesters – any of several polymers with ester linkages in the main chain glycol ethers – via glycol condescension ethoxylates vinyl acetate 1,2-dichloroethane trichloroethylene tetrachloroethylene – also called perchloroethylene; used as a dry cleaning solvent and degreaser vinyl chloride – monomer for polyvinyl chloride polyvinyl chloride (PVC) – a type of plastic used for piping, tubing, other things propylene – used as a monomer and a chemical feedstock isopropyl alcohol – 2-propanol; often used as a solvent or rubbing alcohol acrylonitrile – useful as a monomer in forming Orlon, ABS polypropylene – polymerized propylene propylene oxide polyether polyol – used in the production of polyurethanes propylene glycol – used in engine coolant and aircraft deicer fluid glycol ethers – from the condensation of glycols acrylic acid acrylic polymers allyl chloride epichlorohydrin – chloro-oxirane; used in epoxy resin formation epoxy resins – a type of polymerizing glue from bisphenol A, epichlorohydrin, and some amine butene isomers of butylene – useful as monomers or co-monomers isobutylene – feed for making methyl tert-butyl ether (MTBE) or monomer for copolymerization with a low percentage of isoprene to make butyl rubber 1,3-butadiene (or buta-1,3-diene) – a diene often used as a monomer or co-monomer for polymerization to elastomers such as polybutadiene, styrene-butadiene rubber, or a plastic such as acrylonitrile-butadiene-styrene (ABS) synthetic rubbers – synthetic elastomers made of any one or more of several petrochemical (usually) monomers such as 1,3-butadiene, styrene, isobutylene, isoprene, chloroprene; elastomeric polymers are often made with a high percentage of conjugated diene monomers such as 1,3-butadiene, isoprene, or chloroprene higher olefins polyolefins – such poly-alpha-olefins, which are used as lubricants alpha-olefins – used as monomers, co-monomers, and other chemical precursors. For example, a small amount of 1-hexene can be copolymerized with ethylene into a more flexible form of polyethylene. other higher olefins detergent alcohols Aromatics benzene – the simplest aromatic hydrocarbon ethylbenzene – made from benzene and ethylene styrene – made by dehydrogenation of ethylbenzene; used as a monomer polystyrenes – polymers with styrene as a monomer cumene – isopropylbenzene; a feedstock in the cumene process phenol – hydroxybenzene; often made by the cumene process acetone – dimethyl ketone; also often made by the cumene process bisphenol A – a type of "double" phenol used in polymerization in epoxy resins and making a common type of polycarbonate epoxy resins – a type of polymerizing glue from bisphenol A, epichlorohydrin, and some amine polycarbonate – a plastic polymer made from bisphenol A and phosgene (carbonyl dichloride) solvents – liquids used for dissolving materials; examples often made from petrochemicals include ethanol, isopropyl alcohol, acetone, benzene, toluene, xylenes cyclohexane – a 6-carbon aliphatic cyclic hydrocarbon sometimes used as a non-polar solvent adipic acid – a 6-carbon dicarboxylic acid, which can be a precursor used as a co-monomer together with a diamine to form an alternating copolymer form of nylon. nylons – types of polyamides, some are alternating copolymers formed from copolymerizing dicarboxylic acid or derivatives with diamines caprolactam – a 6-carbon cyclic amide nylons – types of polyamides, some are from polymerizing caprolactam nitrobenzene – can be made by single nitration of benzene aniline – aminobenzene methylene diphenyl diisocyanate (MDI) – used as a co-monomer with diols or polyols to form polyurethanes or with di- or polyamines to form polyureas alkylbenzene – a general type of aromatic hydrocarbon, which can be used as a precursor for a sulfonate surfactant (detergent) detergents – often include surfactants types such as alkylbenzene sulfonates and nonylphenol ethoxylates chlorobenzene toluene – methylbenzene; can be a solvent or precursor for other chemicals benzene toluene diisocyanate (TDI) – used as co-monomers with polyether polyols to form polyurethanes or with di- or polyamines to form polyureas polyurethanes benzoic acid – carboxybenzene caprolactam mixed xylenes – any of three dimethylbenzene isomers, could be a solvent but more often precursor chemicals ortho-xylene – both methyl groups can be oxidized to form (ortho-)phthalic acid phthalic anhydride para-xylene – both methyl groups can be oxidized to form terephthalic acid dimethyl terephthalate – can be copolymerized to form certain polyesters polyesters – although there can be many types, polyethylene terephthalate is made from petrochemical products and is very widely used in petrol stations purified terephthalic acid – often copolymerized to form polyethylene terephthalate polyesters meta-xylene isophthalic acid alkyd resins polyamide resins unsaturated polyesters List of petrochemicals	Technology	Material and chemical	null
67733	https://en.wikipedia.org/wiki/Chestnut	Chestnut	The chestnuts are the deciduous trees and shrubs in the genus Castanea, in the beech family Fagaceae. The name also refers to the edible nuts they produce. They are native to temperate regions of the Northern Hemisphere. Description Chestnut trees are of moderate growth rate (for the Chinese chestnut tree) to fast-growing for American and European species. Their mature heights vary from the smallest species of chinkapins, often shrubby, to the giant of past American forests, C. dentata that could reach . Between these extremes are found the Japanese chestnut (C. crenata) at average; followed by the Chinese chestnut (C. mollissima) at about , then the European chestnut (C. sativa) around . The Chinese and more so the Japanese chestnuts are both often multileadered and wide-spreading, whereas European and especially American species tend to grow very erect when planted among others, with little tapering of their columnar trunks, which are firmly set and massive. When standing on their own, they spread on the sides and develop broad, rounded, dense crowns at maturity. The foliage of the European and American species has striking yellow autumn coloring. Its bark is smooth when young, of a vinous maroon or red-brown color for the American chestnut, grey for the European chestnut. With age, American species' bark becomes grey and darker, thick, and deeply furrowed; the furrows run longitudinally, and tend to twist around the trunk as the tree ages, sometimes reminiscent of a large cable with twisted strands. The leaves are simple, ovate or lanceolate, long and wide, with sharply pointed, widely spaced teeth, with shallow rounded sinuates between. The flowers follow the leaves, appearing in late spring or early summer or into July. They are arranged in long catkins of two kinds, with both kinds being borne on every tree. Some catkins are made of only male flowers, which mature first. Each flower has eight stamens, or 10 to 12 for C. mollissima. The ripe pollen carries a heavy, sweet odor that some people find too sweet or unpleasant. Other catkins have these pollen-bearing flowers, but also carry near the twig from which these spring, small clusters of female or fruit-producing flowers. Two or three flowers together form a four-lobed prickly calybium, which ultimately grows completely together to make the brown hull, or husk, covering the fruits. Chestnut flowers are not self-compatible, so two trees are required for pollination. All Castanea species readily hybridize with each other. Fruit The fruit is contained in a spiny (very sharp) cupule in diameter, also called "bur" or "burr". The burrs are often paired or clustered on the branch and contain one to seven nuts according to the different species, varieties, and cultivars. Around the time the fruits reach maturity, the burrs turn yellow-brown and split open in two or four sections. They can remain on the tree longer than they hold the fruit, but more often achieve complete opening and release the fruits only after having fallen on the ground; opening is partly due to soil humidity. The chestnut fruit has a pointed end with a small tuft at its tip (called "flame" in Italian), and at the other end, a hilum – a pale brown attachment scar. In many varieties, the fruit is flattened on one or two sides. It has two skins. The first one is a hard, shiny, brown outer hull or husk, called the pericarpus; the industry calls this the "peel". Underneath the pericarpus is another, thinner skin, called the pellicle or episperm. The pellicle closely adheres to the seed itself, following the grooves usually present at the surface of the fruit. These grooves are of variable sizes and depths according to the species and variety. The fruit inside these shows a germ with two cotyledons connected to creamy-white flesh throughout. Some varieties have consistently only one embryo per fruit (nut) or have only one large fruit per burr, well rounded (no flat face). The name of varieties with these characteristics may start with "marron" for example marron de Lyon in France, or Marrone di Mugello in Italy. Chestnut fruit may not exhibit epigeal dormancy. It may germinate right upon falling to the ground in the autumn, with the roots emerging from the seed right away and the leaves and stem the following spring. The germ can lose viability soon after ripening and under drying conditions. The superior fruiting varieties among European chestnuts have good size, sweet taste, and easy-to-remove inner skins. American chestnuts are usually very small (around ), but sweet-tasting with easy-to-remove pellicles. Some Japanese varieties have very large nuts (around ), with typically difficult-to-remove pellicles. Chinese chestnut pellicles are usually easy to remove, and their sizes vary greatly according to the varieties, although usually smaller than the Japanese chestnut. Similar species The unrelated horse chestnuts (genus Aesculus) are not true chestnuts, but are named for producing nuts of similar appearance that are mildly poisonous to humans. True chestnuts should also not be confused with water chestnuts, which are tubers of an aquatic herbaceous plant in the sedge family Cyperaceae. Other species commonly mistaken for chestnut trees are the chestnut oak (Quercus prinus) and the American beech (Fagus grandifolia), both of which are also in the Fagaceae family. Brazil nuts, called "Brasil chestnuts" (castañas de Brasil in Spanish) or "chestnuts from Pará" (castanha-do-Pará in Portuguese) are also unrelated. Taxonomy Species Chestnuts belong to the family Fagaceae, which also includes oaks and beeches. The four main species groups are commonly known as American, European, Chinese, and Japanese chestnuts. The taxonomy of the American chestnuts is not completely resolved, particularly between the chinkapins (Castanea ozarkensis and Castanea pumila), which are sometimes considered to be the same species. Genetics have indicated the California native "golden chinkapin" (Chrysolepis chrysophylla) is worthy of inclusion in a different genus along with a species from Coastal China. There is also another chestnut, Castanea alabamensis, which may be its own species. Etymology The name "chestnut" is derived from an earlier English term "chesten nut", which descends from the Old French word chastain (Modern French, châtaigne). The French word in turn derives from Latin Castanea (also the scientific name of the tree), which traces to the Ancient Greek word κάστανον (sweet chestnut). A possible source of the Greek word is the ancient town of Casthanaea in Magnesia. Its location is at the modern village of Keramidi. The town probably took its name, though, from the trees growing around it. In the Mediterranean climate zone, chestnut trees are rarer in Greece because the chalky soil is not conducive to the tree's growth. Kastania is located on one of the relatively few sedimentary or siliceous outcrops. They grow so abundantly there that their presence would have determined the place's name. Still others take the name as coming from the Greek name of Sardis glans (Sardis acorn) – Sardis being the capital of Lydia, Asia Minor, from where the fruit had spread. The name is cited twice in the King James Version of the Bible. In one instance, Jacob puts peeled twigs in the water troughs to promote healthy offspring of his livestock. Although it may indicate another tree, it indicates the fruit was a local staple food in the early 17th century. These synonyms are or have been in use: Fagus Castanea (used by Linnaeus in first edition of Species Plantarum, 1753), Sardian nut, Jupiter's nut, husked nut, and Spanish chestnut (U.S.). Ecology The tree is noted for attracting wildlife. The nuts are an important food for jays, pigeons, wild boar, deer, and squirrels. American and Chinese chinquapins (C. pumila and C. henryi) have very small nuts that are an important source of food for wildlife. Cultivation History Europe and the Near East It has been a staple food in southern Europe, Turkey, and southwestern and eastern Asia for millennia, largely replacing cereals where these would not grow well, if at all, in mountainous Mediterranean areas. Evidence of its cultivation by humans is found since around 2000 BC. Alexander the Great and the Romans planted chestnut trees across Europe while on their various campaigns. A Greek army is said to have survived their retreat from Asia Minor in 401–399 BC thanks to their stores of chestnuts. Ancient Greeks, such as Dioscorides and Galen, wrote of chestnuts to comment on their medicinal properties—and of the flatulence induced by eating too much of it. To the early Christians, chestnuts symbolized chastity. Until the introduction of the potato, whole forest-dwelling communities which had scarce access to wheat flour relied on chestnuts as their main source of carbohydrates. In some parts of Italy, a cake made of chestnuts is used as a substitute for potatoes. In 1583, Charles Estienne and Jean Liébault wrote, "an infinity of people live on nothing else but (the chestnut)". In 1802, an Italian agronomist said of Tuscany that "the fruit of the chestnut tree is practically the sole subsistence of our highlanders", while in 1879 it was said that it almost exclusively fed whole populations for half the year, as "a temporary but complete substitution for cereals". In Britain, boundary records compiled in the reign of King John already showed the famous Tortworth Chestnut in South Gloucestershire, as a landmark; it was also known by the same name of "Great Chestnut of Tortworth" in the days of Stephen. This tree measured over in circumference at from the ground in 1720. The Hundred Horse Chestnut in the chestnut forests on Mount Etna is the oldest living chestnut tree and is said to be even larger. Chestnut trees particularly flourish in the Mediterranean basin. In 1584, the governor of Genoa, which dominated Corsica, ordered all the farmers and landowners to plant four trees yearly, among which was a chestnut tree – plus olive, fig and mulberry trees. Many communities owe their origin and former richness to the ensuing chestnut woods. In France, the marron glacé, a candied chestnut involving 16 different processes in a typically French cooking style, is always served at Christmas and New Year's time. In Modena, Italy, they are soaked in wine before roasting and serving, and are also traditionally eaten on Saint Simon's Day in Tuscany. In the Romagna region, roasted chestnuts are often served with a traditional wine, the Cagnina di Romagna. It is traditional to eat roasted chestnuts in Portugal on St. Martin's Day. Their popularity declined during the last few centuries, partly due to their reputation of "food for poor people". Many people did not want to take chestnut bread as "bread" because chestnut flour does not rise. Some slandered chestnut products in such words as the bread which "gives a sallow complexion" written in 1770, or in 1841 "this kind of mortar which is called a soup". The last decades' worldwide renewal may have profited from the huge reforestation efforts started in the 1930s in the United States to establish varieties of C. sativa which may be resistant to chestnut blight, as well as to relieve the strain on cereal supplies. The main region in Italy for chestnut production is the Mugello region; in 1996, the European Community granted the fruit Protected Geographic Indication (equivalent to the French Appellation d'Origine Contrôlée) status to the Mugello sweet chestnut. It is markedly sweet, peels easily, is not excessively floury or astringent, and has notes of vanilla, hazelnut, and, more subtly, fresh bread. It has no unpleasant aroma, such as yeast, fungus, mold, or paper, which sometimes occur with other chestnuts. The main regions in France for chestnut production are the départements of Ardèche, with the famous "Châtaigne d'Ardèche" (A.O.C) , of the Var (Eastern Provence), of the Cévennes (Gard and Lozère départements) and of the Lyon region. France annually produces over 1,000 metric tons, but still imports about 8,000 metric tons, mainly from Italy. In Portugal's archipelago of Madeira, chestnut liquor is a traditional beverage, and it is gaining popularity with the tourists and in continental Portugal. Asia Always served as part of the New Year's menu in Japan, chestnuts represent both success and hard times—mastery and strength. The Japanese chestnut (kuri) was in cultivation before rice and the Chinese chestnut (C. mollissima) possibly for 2,000 to 6,000 years. During British colonial rule in the mid-1700s to 1947, the sweet chestnut, C. sativa, was widely introduced in the temperate parts of the Indian subcontinent, mainly in the lower to middle Himalayas. They are widely found in British-founded hill stations in northern India, and to a lesser extent in Bhutan and Nepal. They are mainly used as an ornamental tree and are found in almost all British-founded botanical gardens and official governmental compounds (such as larger official residences) in temperate parts of the Indian subcontinent. China has about 300 chestnut cultivars. Moreover, the 'Dandong' chestnut (belonging to the Japanese chestnut C. crenata) is a major cultivar in Liaoning Province. In South Korea, roasted chestnuts (gunbam) are a popular winter snack, and serve as a symbol of abundance in ancestral rituals. Roasted chestnuts are also included in folk songs of Korea, which include "Gunbam Taryeong", a song that celebrates chestnuts, as well as "Jeongseokga", a song from the Goryeo period. Gongju, one of Baekje's former capitals, is renowned for its chestnuts, with an annual chestnut festival that takes place in the winter. In the Samgukji (Records Of The Three Kingdoms), a book that was compiled during the Jin dynasty about the Three Kingdoms, chestnuts are used in the description of Mahan, the former land of Baekje. In the Philippines, the endemic talakatak or Philippine chestnut (Castanopsis philippinensis) is not cultivated commercially, though its nuts are harvested from the wild and consumed locally. Imported chestnuts (known as kastanyas in Tagalog, from Spanish castañas) are traditionally sold as street food in the Philippines during the Christmas season. North America Native Americans were eating the American chestnut species, mainly C. dentata and some others, long before European immigrants introduced their stock to America, and before the arrival of chestnut blight. In some places, such as the Appalachian Mountains, one-quarter of hardwoods were chestnuts. Mature trees often grew straight and branch-free for , up to , averaging up to in diameter. For three centuries, most barns and homes east of the Mississippi River were made from it. In 1911, the food book The Grocer's Encyclopedia noted that a cannery in Holland included in its "vegetables-and-meat" ready-cooked combinations, a "chestnuts and sausages" casserole beside the more classic "beef and onions" and "green peas and veal". This celebrated the chestnut culture that would bring whole villages out in the woods for three weeks each autumn (and keep them busy all winter), and deplored the lack of food diversity in the United States's shop shelves. Soon after that, however, the American chestnuts were nearly wiped out by chestnut blight. The discovery of the blight fungus on some Asian chestnut trees planted on Long Island, New York, was made public in 1904. Within 40 years, the nearly four billion-strong American chestnut population in North America was devastated; only a few clumps of trees remained in Michigan, Wisconsin, California, and the Pacific Northwest. Due to disease, American chestnut wood almost disappeared from the market for decades, although quantities can still be obtained as reclaimed lumber. Today, they only survive as single trees separated from any others (very rare), and as living stumps, or "stools", with only a few growing enough shoots to produce seeds shortly before dying. This is just enough to preserve the genetic material used to engineer an American chestnut tree with the minimal necessary genetic input from any of the disease-immune Asiatic species. Efforts started in the 1930s are still ongoing to repopulate the country with these trees, in Massachusetts and many places elsewhere in the United States. In the 1970s, geneticist Charles Burnham began back-breeding Asian chestnut into American chestnut populations to confer blight resistance with the minimum difference in genes. In the 1950s, the Dunstan chestnut was developed in Greensboro, N.C., and constitutes the majority of blight-free chestnuts produced in the United States annually. Today, the demand for the nut outstrips supply. The United States imported 4,056 metric tons of European in-shell chestnuts worth $10 million in 2007. The U.S. chestnut industry is in its infancy, producing less than 1% of total world production. Since the mid-20th century, most of the US imports are from Southern Italy, with the large, meaty, and richly flavored Sicilian chestnuts being considered among the best quality for bulk sale and supermarket retail. Some imports come from Portugal and France. The next two largest sources of imports are China and South Korea. The French varieties of marrons are highly favored and sold at high prices in gourmet shops. As of 2024, the United States imports 7.5 million pounds of non-organic chestnuts per year. A study of the sector in 2005 found that U.S. producers are mainly part-timers diversifying an existing agricultural business, or hobbyists. Another recent study indicates that investment in a new plantation takes 13 years to break even, at least within the current Australian market. Starting a small-scale operation requires a relatively low initial investment; this is a factor in the small size of the present production operations, with half of them being between . Another determining factor in the small productivity of the sector is that most orchards have been created less than 10 years ago, so have young trees which are as now barely entering commercial production. Assuming a yield for a 10-year-old tree is a reliable conservative estimate, though some exceptional specimens of that age have yielded . So, most producers earn less than $5,000 per year, with a third of them not having sold anything so far. Moreover, the plantings have so far been mostly of Chinese species, but the products are not readily available. The American Chestnut Foundation in collaboration with many partners (SUNY ESF, the American Chestnut Cooperators' Foundation and many others from education, research, and industry sectors contributing to the program) are in the last stages of developing a variety that is as close as possible to the American chestnut, while having incorporated the blight-resistant gene of the Asiatic species. Considering the additional advantage that chestnut trees can be easily grown organically, and assuming the development of brands in the market and everything else being equal, home-grown products would reach higher prices than imports, the high volume of which indicates a market with expanding prospects. As of 2008, the price for chestnuts sold fresh in the shell ranges from $1.50/lb ($3.30/kg) wholesale to about $5/lb ($11/kg) retail, depending mainly on the size. Australia and New Zealand The Australian gold rush of the 1850s and 1860s led to the first recorded plantings of European chestnut trees, brought from Europe by settlers. Along the years, most chestnut tree plantations were C. sativa stock, which is still the dominant species. Some of these remain today. Some trees in northern Victoria are around 120 years old and up to 60 m tall. Chestnuts grow well in southwest Western Australia, which has cold winters and warm to hot summers. As of 2008, the country has nearly 350 growers, annually producing around 1,200 metric tons of chestnuts, of which 80% come from northeast Victoria. The produce is mostly sold to the domestic fresh fruit market. Chestnuts are slowly gaining popularity in Australia. A considerable increase in production is expected in the next 10 years, due to the increase in commercial plantings during the last 15 to 25 years. By far, the most common species in Australia is the European chestnut, but small numbers of the other species, as well as some hybrids, have been planted. The Japanese chestnut (C. crenata) does well in wet and humid weather and in hot summers (about 30 °C); and was introduced to New Zealand in the early 1900s, more so in the upper North Island region. Cultivation ecology Climate and seasonal germination cycle Chestnuts produce a better crop when subjected to chill temperatures during the dormant period. Frosts and snowfalls are beneficial rather than harmful to the trees. The dormant plant is very cold-hardy in Britain, to the Royal Horticultural Society's H6 hardiness rating, to -20 °C. Chestnut is hardy to USDA zone 5, which is lower in average minimal temperature than London in zone 9. The young growth in spring, even on mature plants, however, is frost-tender; bud-burst is later than most other fruit trees, so late frosts can be damaging to young buds. Trees can be found at altitudes between 200 and above sea level; some mention between 300 and altitude, while the famous Hundred Horse Chestnut on Mount Etna stands at 1200 metres. They can tolerate maritime exposure, although growth is reduced. Seeds germinate in late winter or early spring, but the life length is short. If kept moist, they can be stored in a cool place for a few months, but must be checked regularly for signs of germination. Low temperature prolongs dormancy. Sowing them as soon as ripe is better, either in cold frames or seedbeds outdoors, where they can be left in situ for one to two years before being planted in their permanent positions, or in pots, where the plants can be put out into their permanent positions in summer or autumn. They must be protected from the cold in their first winter, and also from mice and squirrels. Chestnuts are considered self-sterile, so at least two trees are needed for pollination. Soil requirements Castanea grows best in a soil with good drainage and adequate moisture. The tree prefers sloping, deep soils; it does not like shallow or heavy soils with impermeable, clay subsoils. The Chinese chestnut prefers a fertile, well-drained soil, but it grows well in fairly dry, rocky, or poor soils. Although Castanea can grow in very acidic soil, and while these soils are reasonably well tolerated, the preferred range is from pH 5.5-6.0. It does not grow well on alkaline soils, such as chalk, but thrives on soils such as those derived from granite, sandstone, or schist. On alkaline soils, chestnut trees can be grown by grafting them onto oak rootstocks. Recently cleared land is best avoided to help resist the root rot, Armillaria mellia. Sun exposure Castanea likes a full sun position. An experiment with C. dentata seedlings in Ohio confirmed the need for sun for optimal growth. The butt of the tree is sometimes painted with white paint to protect the tree from sunburn until it has developed enough canopy. Wide spacing between the trees encourages low, broad crowns with maximum exposure to sunshine to increase fruit production. Where chestnut trees touch, virtually no fruit is produced. Current industrial planting spacings can range from 7 x 7 to 20 x . The closer plantings, which are more popular, mean quicker increases in short-term production, but heavy pruning or even tree removal is required later. Watering The optimum rainfall for chestnut trees is or more, ideally in even distribution throughout the year. Mulching during summer is recommended. Rainfall below per year needs be complemented with, for example, a drip irrigation system. This should water the soil at the outer half of the circle formed by the drip line to encourage root growth. Independently from annual rainfall, watering young trees is recommended at least during summer and early autumn. Once established, they resist droughts well. Preservation In addition to being consumed fresh, chestnuts can also be canned, pureed, or preserved in sugar or syrup (marrons glacés). Shelled and cooked nuts should be covered, refrigerated, and used within 3–4 days. Cooked chestnuts, either whole, chopped, or pureed, may be frozen in an airtight container and held up to 9 months. Because of their high water content, transpiration rates, and consequent loss weight, the nuts react as fresh fruits (not as nuts). They should be kept cool at all times, including in shops when on display for sale. To preserve their freshness for a few months with no artificial refrigeration, the chestnuts can be soaked in cold water for about 20 hours immediately after harvest, after which they are dried in the shade, then layered in dry sand. Chestnuts behave similarly to seeds in that they produce very little ethylene, and their respiration rate is low, varying between 5 and 20 mg/(kg·h) depending on the temperature. Pests Mammals and birds Grey squirrels strip bark from when the tree is about eight years old and onward through the life of the tree. Rabbits and wallabies can do great damage to young trees, which need guarding by some fence or by wrapping the tree trunk in sisal or other appropriate material. Deer and kangaroos can also be troublesome. Cattle and horses may require temporary fencing to prevent them from damaging fallen chestnuts at harvest time. The sulphur-crested cockatoo can damage branches up to in diameter by carrying out "beak maintenance" on young trees. Rosellas can be troublesome at harvest time. Shrews, squirrels, mice, and other critters often eat the chestnut seed after it has sprouted within the first, and even second years of growth. Some avoid this by removing the chestnut seed from the stem. Insects Dryocosmus kuriphilus, the oriental chestnut gall wasp, is native to China, but is an invasive pest elsewhere. It attacks and destroys the chestnut fruit. It is considered the world's worst pest of chestnuts. The larvae of the polyfag moth (Phytomyza horticola) species are among those that do most damage to shoots and foliage. The most frequently occurring pests are the winter moth (Operophtera brumata) and the mottled umber moth (Erannis defoliaria). The oak roller weevil (Attelabus nitens) causes relatively less damage by rolling up the leaves into a barrel shape to shelter its eggs and developing larvae. The insects swarm from the end of April to mid-June, and damage the tree's flower buds during their feeding season. The larvae of the oak-leaf-mining moth, also called the tischerid moth (Tischeria ekebladella), digs white, see-through mines in chestnut leaves. It lays its eggs in the leaves between May and June. The larvae cause white spots in the leaves by chewing them from the inside. The oak aphid (Myzocallis castanicola) sucks on the apex of young shoots and leaves. Native to Europe and North America, it is, for example, active in Hungary. Leaves do not roll up, but their feeding delays the growth of shoots and damages young graft-shoot hosts. Commercial plantations and nurseries spray pesticides during the shoots' growth period to fight the damage. The chestnut mosaic virus is probably transmitted by M. castanicola aphids. The chestnut weevil (Curculio elephas) most often damages the fruits. In Hungary, it swarms in chestnut orchards around August 20, particularly strongly around noon and in sunny weather. The eggs are laid into the cupules or around the peduncle joints. The larvae feed on the nuts and leave only nutchips and excrement within. While the chestnuts ripen, the larvae retreat into the ground after having chewed their way out of the nuts. The following July, they turn into pupae. The larvae of the chestnut weevil can only chew their way out of a fallen nut, so breeding occurs mostly where chestnuts lie on the ground for a sufficient length of time, or where the trees produce many small fruits which remain behind at the harvest. Timing the harvests to pick up the chestnuts as soon as they fall reduces the numbers of the overwintering larvae. Regular soil work is also unfavourable to its life habits. Small grafts are sprayed with chemicals. A warm, aerosol-based protection has been developed for older trees, by Sifter and Bürgés in 1971. Planting chestnut orchards beside turkey oak forests is not advised, because both trees are susceptible to the chestnut weevil (which also uses the turkey oak acorn to develop), and the turkey oak trees can pass it on to the chestnut trees. In Hungary, the most common moth threatening chestnut trees is the acorn moth (Laspeyreisa splendana) and its subspecies. Its grayish-yellow larvae cause similar damage to that of the chestnut weevil, but they spin characteristic webs among the nutchips and larval excrement. This moth causes about 5–41% of the damage that occurs in western Hungary's plantations. Plantations need regular protection against these moths, the occurrence of which does not decrease. In New Zealand, the grass grub beetle eats the soft, new-season foliage. They can entirely strip a young tree in the late spring, when they fly at dusk, often in huge numbers. Diseases Chestnut blight fungus (Cryphonectria parasitica) (formerly Endothia parasitica) affects chestnut trees. The Eastern Asian species have coevolved with this disease and are moderately to very resistant to it, while the European and North American species, not having been exposed to it in the past, have little or no resistance. Early in the 20th century, chestnut blight destroyed about four billion American chestnut trees, and reduced the most important tree throughout the East Coast to an insignificant presence. The American chinkapins are also very susceptible to chestnut blight. The European and West Asian chestnuts are susceptible, but less so than the American species. The resistant species (particularly Japanese and Chinese chestnut, but also Seguin's chestnut and Henry's chestnut) have been used in breeding programs in the U.S. to create hybrids with the American chestnut that are also disease-resistant. The bark miner Spulerina simploniella (Lepidoptera: Gracilariidae) was found in intensively managed chestnut coppices in Greece, but not in orchards. The larvae (and the rain) may be agents in the spread of the disease. They mine under the thin periderm of young trees up to 10 years old, while the stem bark is still smooth. Rain during the pupation period (around the last week of May and first two weeks of June), and the actions of the larvae, may collude for conidiospores to come into contact with the freshly exposed phloem, thus causing cankers. Ink disease also appears in a number of other plants. The disease attacks the phloem tissue and the cambium of the roots and root collars about 10–20 cm above ground. Wet rot settles in as a result. It was named after the ink-black color of the tannic acid becoming (oxidized) after seeping out, but that symptom is not a characteristic of only that disease. The same ink-black color can appear following other types of decays and mechanical injuries that make liquids seep through; these liquids can also oxidize after contact with air. Moreover, with some phytophthoric diseases, no tannic acid is generated. With the ink disease, the leaves turn yellow and later fall off; the fruits remain small, and the nuts prematurely drop out of the burrs. These dry and remain on the trees throughout winter. In acute cases, root decay makes the trees dry out and wither away. It is caused by Phytophthora cambivora and Phytophthora cinnamomi. Phytophthora disease is the longest-known chestnut tree disease leading to tree death. Of the two main pathogens for this disease, the one in European chestnuts is known since 1971 to be Phytophthora cambivora. Phytophthora cinnamomi was discovered in chestnut trees in the United States in 1932. Both trigger similar symptoms. Since then, it has also been shown to occur in most European chestnut-growing countries. Differentiating between the two pathogens is difficult. Chemicals seem of little effectiveness. Many countries impose strict prophylactic rules to prevent the spread of the disease. Melanconis modonia can infect trees through injuries and induce "bark death". It was first reported in Hungary by Hausz in 1972. The damage is of little consequence in older or stronger trees, but it affects sapling graftings in nurseries. Coryneum perniciosum, one of the two conidium-like side forms of this fungus, occurs on all decayed, ligneous parts of a chestnut tree. The symptoms of infection on young, smooth trunks is similar to that of the chestnut blight fungus Cryphonectria. For this reason, it has persistently been wrongly thought of as the pathogen for ink disease. With Melanconis, the bark sinks in and takes on brownish-red tones, with black, lentil-like multicell conidium bodies and black cone-like stromata breaking through the bark. Unlike with Cryphonectria, though, no orange-colored fruiting bodies are seen. Prevention primarily includes keeping trees in good shape; some further protections against Cryphonectria also help prevent bark death caused by Melanconis. Chestnut mosaic virus is probably transmitted by the oak aphid Myzocallis castanicola. Root rot is caused by the honey fungus Armillaria mellia. When planting Castanea, recently cleared land is best avoided to help resist this fungus. The disease is more prevalent on heavier and poorly drained soil types. Leaf spot is the most common disease for chestnut trees (Mycosphaerella maculiformis). It is known as cylindrosporium leaf spot disease, after its summer conidium form Cylindrosporium castaneae. The pathogens spend the winter in the white spots of the fallen leaves. At spring time, it reinfects the new leaves. In or near June, tiny white spots on the leaves appear, which grow and turn brown over time. At the end of the summer, the spots entirely cover the leaf, which turns yellow. In rainy and humid weather with large temperature fluctuation, the tree loses its leaves. If August is dry and warm, the infected leaves roll up, the arteries twist, and the dead leaves dry on the tree until defoliage. This recurs yearly, though the extent of the damage varies from year to year. Some species are more resistant than others. Oak mildew is among several foliage diseases of smaller significance for European chestnut growing. It infects the most trees (Microsphaera alphitoides). Younger trees suffer most; their shoots become short-jointed, growth is delayed, and they develop sensitivity to frostbite. In older trees, the fungus usually infects only the tip of the shoots. The pathogens hibernate in the shoots and infect the leaves from there. The fungus grows on the top of the leaves, with the appearance of a coating only in midsummer. The infected leaves' development slows down or stops, the distance between their vessels shrinks, and the vessels themselves become curly. In storage rot, breaking the tuft provides the most common entrance for fungal spores during storage. Ciboria, the most diffuse, turns the flesh black and spongy. Other fungi are known, such as Rhizopus, Fusarium, and Colletotrichum. In chestnuts, Colletotrichum disease symptoms may also be called blossom end rot. Browning of the chestnut burs at the blossom end may be a first sign in August. At harvest time, blackening of pointed end of the chestnut shell and kernel indicates infection. The extent of blackening can vary. It can range from a barely visible black tip of the kernel to the whole nut being black. Parts of the nut kernel with no color change remain edible. Chestnut canker can be caused by fungi of genus Dendrostoma. Coppicing Most chestnut wood production is done by coppice systems, cut on a 12-year rotation to provide small timber which does not split as badly as large logs. In southern England (particularly in Kent), sweet chestnut has traditionally been grown as coppices, being recut every 10 years or so on rotation for poles used for firewood, and fencing (fence posts and chestnut paling). Sustainable forest management An excellent soil-enriching understory in pine forests, sustainable forest management incorporates more mixed plantings of proven efficiency, as opposed to monosylviculture. A study presented in 1997 has evaluated positively the potential increase in productivity with mixed stands and plantations, compared to plots of only one species. The relative yield total values of the mixed plantings steadily increase with time. C. sativa responds well to competitive pressure from Pseudotsuga menziesii, the latter also showing a higher productivity. C. dentata seedlings in Ohio reforestation efforts are best achieved by planting them in places with little or no arboreous land cover, because of the need for light. Production In 2020, world production of cultivated chestnuts was 2,322 tonnes, led by China with 75% of the total (table). Uses Nutrition Chestnuts depart from the norm for culinary nuts, as they have little protein or fat; their calories come chiefly from carbohydrates. Fresh chestnut fruits provide about of food energy per 100 g of edible parts, which is much lower than walnuts, almonds, other nuts, and dried fruit (about per 100 g). In some areas, sweet chestnut trees are called "bread trees". When chestnuts are just starting to ripen, the fruits are mostly starch and are firm under finger pressure from the high water content. As the chestnuts ripen, the starch is slowly converted into sugars, and moisture content decreases. Upon pressing the ripe chestnut, a slight "give" can be felt; the hull is not so tense, and space occurs between the flesh of the fruit and it. Raw chestnuts are 60% water and contain 44 grams of carbohydrates, 2 grams of protein, one gram of fat, supplying 200 calories in a 100-gram reference amount (table). Chestnuts provide some B vitamins and dietary minerals in significant content (table). Their carbohydrate content compares with that of wheat and rice. Chestnuts have twice as much starch as the potato on an as-is basis. They contain about 8% of various sugars, mainly sucrose, glucose, fructose, and in lesser amounts, stachyose and raffinose, which are fermented in the lower gut, producing gas. Chestnuts are among the few "nuts" that contain vitamin C, with 48% of the Daily Value in a 100-gram serving (table). The amount of vitamin C decreases by roughly 40% upon heating (typically, the vitamin is decreased or destroyed in heated foods). Fresh chestnuts contain about 52% water by weight, which evaporates relatively quickly during storage. They can lose as much as 1% of weight in one day at 20 °C (68 °F) and 70% relative humidity. Culinary The fruit can be peeled and eaten raw, but it can be somewhat astringent, especially if the pellicle is not removed. Another method of eating the fruit involves roasting, which does not require peeling. Roasting requires scoring the fruit beforehand to prevent explosion of the fruit due to expansion. Once cooked, its texture is slightly similar to that of a baked potato, with a delicate, sweet, and nutty flavour. This method of preparation is popular in many countries, where the scored chestnuts may be cooked mixed with a little sugar. Chestnuts can be dried and milled into flour, which can then be used to prepare breads, cakes, pies, pancakes, pastas, polenta (known in Corsica as pulenda), or used as thickener for stews, soups, and sauces. Chestnut cake may be prepared using chestnut flour. In Corsica, the flour is fried into doughnut-like fritters called fritelli and made into necci, pattoni, , and cialdi. The flour can be light beige like that from Castagniccia, or darker in other regions. It is a good solution for long storage of a nutritious food. Chestnut bread can stay fresh as long as two weeks. The nuts can also be eaten candied, boiled, steamed, deep-fried, grilled, or roasted in sweet or savory recipes. They can be used to stuff vegetables, poultry, fowl, and other edibles. They are available fresh, dried, ground, or canned (whole or in puree). Candied chestnuts (whole chestnuts candied in sugar syrup, then iced) are sold under the French name marrons glacés or Turkish name kestane şekeri ("sugared chestnuts"). They appeared in France in the sixteenth century. Toward the end of nineteenth century, Lyon went into a recession with the collapse of the textile market, notably silk. Clément Faugier, a civil engineer, was looking for a way to revitalize the regional economy. In 1882 at Privas, he invented the technology to make marrons glacés on an industrial scale (although a great number of the more than 20 necessary steps from harvest to the finished product are still accomplished manually). Chestnuts are picked in autumn, and candied from the start of the following summer for the ensuing Christmas. Thus, the marrons glacés eaten at Christmas are those picked the year before. In Spain, on 31 October on the eve of the All Saints' Day, Catalonia celebrates la castanyada, a festivity that consists of eating chestnuts, panellets, sweet potatoes and muscatell. On November, in the regions of Galicia, Asturias, Cantabria and other Northern provinces and Portugal, the Magosto is celebrated. In Hungarian cuisine, cooked chestnuts are puréed, mixed with sugar (and usually rum), forced through a ricer, and topped with whipped cream to make a dessert called gesztenyepüré (chestnut purée). In Swiss cuisine, a similar dish made with kirsch and butter is called vermicelles. A French version is known as "Mont Blanc". A fine granular sugar can be obtained from the fermentation of the juice, as well as a beer; the roasted fruit provides a coffee substitute. Parmentier, who among other things was a famous potato promoter, extracted sugar from chestnuts and sent a chestnut sugarloaf weighing several pounds to the Academy of Lyon. The continental blockade following shortly after (1806–1814) increased the research into developing chestnuts as a source of sugar, but Napoleon chose beets instead. Sweet chestnuts are not easy to peel when cold. One kilogram of untainted chestnuts yields about 700 g of shelled chestnuts. Animal fodder and litter Chestnuts are often added to animal fodder. A first soak in limewater removes their bitter flavour, then they are ground and mixed with the ordinary provender. Other methods of preparation are also used. It is given to horses and cattle in the Orient, and to pigs in England, France and other places. The leaves are not as prone to be insect-eaten as those of the oak, and are also used for fodder. Timber Chestnut is of the same family as oak, and likewise its wood contains many tannins. This renders the wood very durable, gives it excellent natural outdoor resistance, and saves the need for other protection treatment. It also corrodes iron slowly, although copper, brass, or stainless metals are not affected. Chestnut timber is decorative. Light brown in color, it is sometimes confused with oak wood. The two woods' textures are similar. When in a growing stage, with very little sap wood, a chestnut tree contains more timber of a durable quality than an oak of the same dimensions. Young chestnut wood has proved more durable than oak for woodwork that has to be partly in the ground, such as stakes and fences. After most growth is achieved, older chestnut timber tends to split and warp when harvested. The timber becomes neither so hard nor so strong as oak. The American chestnut C. dentata served as an important source of lumber, because it has long, unbranched trunks. In Britain, chestnut was formerly used indiscriminately with oak for the construction of houses, millwork, and household furniture. It grows so freely in Britain that it was long considered a truly native species, partly because the roof of Westminster Hall and the Parliament House of Edinburgh were mistakenly thought to be constructed of chestnut wood. Chestnut wood, however, loses much of its durability when the tree is more than 50 years old, and despite the local chestnut's quick growth rate, the timber used for these two buildings is considerably larger than a 50-year-old chestnut's girth. It has been proven that the roofs of these buildings are made of Durmast oak, which closely resembles chestnut in grain and color. It is therefore uncommon to find large pieces of chestnut in building structures, but it has always been highly valued for small outdoor furniture pieces, fencing, cladding (shingles) for covering buildings, and pit-props, for which durability is an important factor. In Italy, chestnut is also used to make barrels used for aging balsamic vinegar and some alcoholic beverages, such as whisky or lambic beer. Of note, the famous 18th-century "berles" in the French Cévennes are cupboards cut directly from the hollowed trunk. Fuel Dry, chestnut firewood is best burned in a closed log-burner, because of its tendency to spit when on an open fire. Leather Chestnut wood is a useful source of natural tannin and was used for tanning leather before the introduction of synthetic tannins. On a 10% moisture basis, the bark contains 6.8% tannin and the wood 13.4%. The bark imparts a dark color to the tannin, and has a higher sugar content, which increases the percentage of soluble non-tans, or impurities, in the extract; so it was not employed in this use. Chestnut tannin is obtained by hot-water extraction of chipped wood. It is an ellagic tannin and its main constituents are identified by castalagin (14.2%) and vescalagin (16.2%).It has a naturally low pH value, relatively low salts content, and high acids content. This determines its astringency and its capability to fix raw hides. These properties make chestnut extract especially suitable for the tanning of heavy hides and to produce leather soles for high-quality shoes in particular. It is possible to obtain a leather with high yield in weight, which is compact, firm, flexible, and waterproof. Chestnut-tanned leathers are elastic, lightfast, resistant to traction and abrasion, and have warm color. Chestnut tannin is one of the pyrogallol class of tannins (also known as hydrolysable tannin). As it tends to give a brownish tone to the leather, it is most often used in combination with quebracho, mimosa, tara, myrabolans, and valonia. The wood seems to reach its highest tannin content after the trees reach 30 years old. The southern European chestnut wood usually contains at least 10 to 13% more tannin than chestnut trees in northern climates. Other uses Fabric can be starched with chestnut meal. Linen cloth can be whitened with chestnut meal. The leaves and the skins (husk and pellicle) of the fruits provide a hair shampoo. Hydrolysable chestnut tannins can be used for partial phenol substitution in phenolic resin adhesives production and also for direct use as resin. Chestnut buds have been listed as one of the 38 substances used to prepare Bach flower remedies, a kind of alternative medicine promoted for its effect on health. However, according to Cancer Research UK, "there is no scientific evidence to prove that flower remedies can control, cure or prevent any type of disease, including cancer". In culture In the film based on the novel by E. M. Forster, Howards End, Mrs. Ruth Wilcox (Vanessa Redgrave) tells of her childhood home, where superstitious farmers would place pigs' teeth in the bark of the chestnut trees and then chew on this bark to ease toothaches. In the novel, the tree is actually a Wych elm. Under the Spreading Chestnut Tree is a set of variations, with fugue, for orchestra composed in 1939 by Jaromír Weinberger. In Honoré de Balzac's novel Père Goriot, Vautrin states that Eugène de Rastignac's family is living off chestnuts; this symbolism is used to represent how impoverished Eugene's family is. "The Christmas Song" famously mentions chestnuts in its opening line, and is commonly subtitled "Chestnuts Roasting on an Open Fire." Notable specimens Hundred Horse Chestnut on Mount Etna, 57.9 m (190 ft) circumference in 1780, (64-meter circumference in 1883) Tortworth Chestnut. 15.8-meter (52 ft) circumference in 1776, when it was described as "the largest tree in England" Sacred Chestnut of Istán, circumference, estimated to be between 800 and 1,000 years old.	Biology and health sciences	Fagales	null
67740	https://en.wikipedia.org/wiki/American%20chestnut	American chestnut	The American chestnut (Castanea dentata) is a large, fast-growing deciduous tree of the beech family native to eastern North America. As is true of all species in the genus Castanea, the American chestnut produces burred fruit with edible nuts. The American chestnut was once common in its Appalachian Mountain range and was a dominant species in the oak-chestnut forest region of its central and southern range. During the early to mid-20th century, American chestnut trees were devastated by chestnut blight, a fungal disease that came from Japanese chestnut trees that were introduced into North America from Japan. It is estimated that the blight killed between three and four billion American chestnut trees in the first half of the 20th century, beginning in 1904. Few mature American chestnuts exist within its former range, although many stumps and root systems continue to send up saplings. Most of these saplings get infected by chestnut blight, which girdles and kills them before they attain maturity. There are hundreds of large ( in diameter) American chestnuts outside its historical range, some in areas where less virulent strains of the pathogen are more common, such as the 600 to 800 large trees in Northern Michigan. The species is listed as endangered in Canada under the Species at Risk Act. American chestnuts are also susceptible to ink disease, particularly in the southern part of its native range; this likely contributed to the devastation of the species. Several groups are attempting to create blight-resistant American chestnuts. Scientists at the SUNY College of Environmental Science and Forestry created the Darling 58 cultivar of American chestnut by inserting the oxalate oxidase gene from wheat into the genome of an American chestnut. When expressed in the vascular cambium of the Darling 58 cultivar, the oxalate oxidase enzyme degrades the oxalic acid produced by the chestnut blight, reducing damage to the vascular cambium and resisting girdling of the trunk. As of 2021, the researchers who developed this cultivar are working toward applying for government permission to make these trees available to the public. If approved, these chestnut trees would be the first genetically modified forest trees released into the wild in the United States. Alternate approaches to developing a blight-resistant cultivar include cross-breeding among partially blight-resistant American chestnuts or crossbreeding with the moderately blight-resistant Chinese chestnut, then backcrossing with the American chestnut, with the goal of retaining most of its genes. Description Castanea dentata is a rapidly-growing, large, deciduous hardwood eudicot tree. A singular specimen manifest in Maine has attained a height of Pre-blight sources give a maximum height of , and a maximum circumference of . Post-blight sources erroneously report a greater maximum size of the species compared with pre-blight, likely due to nostalgia, to interpretations of pre-blight measurements of circumference as being measurements of diameter, and to the misapprehension that pre-blight observations of maximum size represented observations of average size. It is considerably larger than the closely related Allegheny chinquapin (Castanea pumila). There are several other chestnut species, such as the European sweet chestnut (C. sativa), Chinese chestnut (C. mollissima), and Japanese chestnut (C. crenata). Castanea dentata can be distinguished by a few morphological traits, such as petiole length, nut size and number of nuts per burr, leaf shape, and leaf size, with leaves being long and broad—slightly shorter and broader than the sweet chestnut. It has larger and more widely spaced saw-teeth on the edges of its leaves, as indicated by the scientific name dentata, Latin for "toothed". The European sweet chestnut was introduced in the United States by Thomas Jefferson in 1773. The European sweet chestnut has hairy twig tips in contrast to the hairless twigs of the American chestnut. This species has been the chief source of commercial chestnuts in the United States. Japanese chestnut was inadvertently introduced into the United States by Thomas Hogg in 1876 and planted on the property of S. B. Parsons in Flushing, New York. The Japanese chestnut has narrow leaves, smaller than either American chestnut or sweet chestnut, with small, sharply-pointed teeth and many hairs on the underside of the leaf and is the most blight-resistant species. The chestnut is monoecious, and usually protandrous producing many small, pale green (nearly white) male flowers found tightly occurring along 6 to 8 inch long catkins. The female parts are found near the base of the catkins (near twig) and appear in late spring to early summer. Like all members of the family Fagaceae, American chestnut is self-incompatible and requires two trees for pollination, which can be with other members of the Castanea genus. The pollen of the American chestnut is considered a mild allergen. The American chestnut is a prolific bearer of nuts, with inflorescence and nut production in the wild beginning when a tree is 8 to 10 years old. American chestnut burrs often open while still attached to the tree, around the time of the first frost in autumn, with the nuts then falling to the ground. American chestnut typically have three nuts enclosed in a spiny, green burr, each lined in a tan velvet. In contrast, the Allegheny chinquapin produces but one nut per burr. Evolution and ecology Chestnuts are in the Fagaceae family along with beech and oak. Chestnuts are not closely related to the horse chestnut, which is in the family Sapindaceae. Phylogenetic analysis indicates a westward migration of extant Castanea species from Asia to Europe to North America, with the American chestnut more closely related to the Allegheny chinquapin (Castanea pumila v. pumila) than to European or Asian clades. The genomic range of chestnuts can be roughly divided into a clinal pattern of northeast, central, and southwest populations, with southwest populations showing greatest diversity, reflecting an evolutionary bottleneck likely due to Quaternary glaciation. Two lineages of American chestnut have been identified, one a hybrid between the American chestnut and the Allegheny chinquapin from the southern Appalachians. The other lineage of American Chestnut shows a gradual loss of genetic diversity along a Northward vector, indicating possible expansion of range following the most recent Glacial Maximum during the Wisconsin glaciation. Ozark chinkapin, which is typically considered either a distinct species (C. ozarkensis) or a subspecies of the Allegheny chinquapin (C. pumila subsp. ozarkensis), may be ancestral to both the American chestnut and the Allegheny chinquapin. A natural hybrid of C. dentata and C. pumila has been named Castanea × neglecta. The American chestnut population was reduced to 1–10% of its original size as a result of the chestnut blight and has not recovered. The surviving trees are "frozen in time" with shoots re-sprouting from survivor rootstock but almost entirely undergoing blight-induced dieback without producing chestnuts. Unexpectedly, American chestnut appears to have retained substantial genetic diversity following the population bottleneck, which is at odds with the limited incidence of blight resistance/tolerance in extant populations. The pre-blight distribution of the American chestnut was restricted to moist, but well-drained, steep slopes with acid loam soils. According to analysis of old forest dust data, the tree was rare or absent in New England prior to 2,500 years before the present, but rapidly established dominance in these forests, becoming a common tree over a range from Maine and southern Ontario to Mississippi, and from the Atlantic coast to the Appalachian Mountains and the Ohio Valley. Within its range, the American chestnut was the dominant timber of mountain ridges and sandstone soils. Along the Blue Ridge Mountains of North Carolina, the American chestnut dominated the area above the range of the Eastern hemlock and below 1,500 meters. In Western Maryland, the American chestnut comprised 50% of ridge timber and 36% of forested slopes. The tree's abundance was due to a combination of rapid growth, relative fire resistance, and a large annual nut crop, in comparison to oaks, which do not reliably produce sizable numbers of acorns every year. Fire was common in the pre-blight ecosystem of the American chestnut, perhaps in part due to unique traits of the tree, including fire tolerance, highly flammable litter, tall stature, rapid growth, and ability to resprout. Historically, the mean fire return interval was 20 years or less in chestnut-predominant ecologies, with a forest stand pattern that was more open than is currently the case. The American chestnut was an important tree for wildlife, providing much of the fall mast for species such as white-tailed deer, wild turkey, Allegheny woodrat and (prior to its extinction) the passenger pigeon. Black bears were also known to eat the nuts to fatten up for the winter. The American chestnut also contains more nitrogen, phosphorus, potassium and magnesium in its leaves than other trees that share its habitat, so they return more nutrients to the soil which helps with the growth of other plants, animals, and microorganisms. The American chestnut is preferred by some avian seed hoarders, and was particularly important as a food source during years where the oak mast failed. The functional extinction of the American chestnut may have resulted in the extinction of some of the tree's host-specialist insect associates, including the Greater Chestnut Weevil. Parasites of American chestnut The appearance of invasive pathogens of the American Chestnut into the eastern deciduous forest ecosystem is just one instance of the Columbian exchange of pathogens. While the Columbian exchange moved valuable crops between the Americas, Europe and Asia, there was also a downside, as the rapid introduction of invasive and unfamiliar pathogens resulted in serious damage or extinction of some host species. Chestnut blight Prior to the Chestnut blight, the American chestnut was a dominant tree in the ecosystem of the eastern deciduous forest. It was said that a squirrel could walk from New England to Georgia solely on the branches of American chestnuts. Once an important hardwood timber tree, the American chestnut suffered a catastrophic population collapse due to the chestnut blight, a disease caused by an Asian bark fungus (Cryphonectria parasitica, formerly Endothia parasitica). The fungus was introduced when infected Japanese chestnut trees were brought to North America in the late 19th century. Chestnut blight was first noticed on American chestnut trees in what was then the New York Zoological Park, now known as the Bronx Zoo, in the borough of The Bronx, New York City, in 1904, by chief forester Hermann Merkel. Merkel estimated that, by 1906, blight had infected 98 percent of the chestnut trees in the borough. While Asian chestnut species evolved with the blight and developed a strong resistance, the American chestnut and Allegheny chinquapin have little resistance. The airborne bark fungus spread a year and in a few decades girdled and killed more than three billion American chestnut trees. Salvage logging during the early years of the blight may have unwittingly destroyed trees that had high levels of resistance to the disease and thus aggravated the calamity. New shoots often sprout from the roots when the main stem dies, so the species has not yet become extinct. However, the stump sprouts rarely reach more than in height before blight infection returns, so the species is classified as functionally extinct since the chestnut blight only actively kills the above ground portion of the American chestnut tree, leaving behind the below-ground components such as the root systems. It was recorded in the 1900s that the chestnut blight would commonly reinfect any novel stems that grew from the stumps, therefore maintaining a cycle that would prevent the American chestnut tree from re-establishing. However, some American chestnut trees have survived because of a small natural resistance to the chestnut blight. The high density of American chestnuts within its range and the lack of natural immunity allowed the blight to spread quickly and cause infection and die-off in nearly every tree exposed. Early attempts to treat chestnut blight were both chemical, such as the use of fungicides, and physical, such as removing infected limbs through tree surgery and the removal of infected trees from cultivated and wild stands. Quarantine measures were also put into place, with the later support of the Plant Quarantine Act, which was an attempt to prevent the importation of other potential plant pathogens. These attempts to contain the spread of chestnut blight were unsuccessful; the devastation of the species was worsened because the chestnut blight resulted in isolation of remaining specimens, resulting in asexual propagation of many isolated American chestnuts, low genetic diversity of stands of American chestnuts and consequent vulnerability to extirpation. Chestnut blight is not to be confused with sun scald, where winter sun reflects off of snow, warming the bark on the sun-facing trunk (this is the south-facing trunk in the Northern Hemisphere). This snow-reflected sunlight repeatedly warms and thaws the trunk during the day, resulting in vulnerability of the bark and cambium to freezing cold temperatures during the subsequent night, eventually resulting in bark cankers that resemble chestnut blight. Also, sun scald makes the damaged bark vulnerable to invasion by pathogens. Ink disease Before the onset of chestnut blight and prior to 1824, an epidemic of ink disease struck American chestnuts, most likely brought to the southern United States on Cork oak trees imported from Portugal. This fungal pathogen is known to also kill the roots and collars of several Castanea species, including the European species sweet chestnut (Castanea sativa). It affected primarily chestnuts in the Southeastern US and at the later time when chestnut blight struck, the range of C. dentata may have already been reduced. The potential range of blight-resistant American chestnuts is substantially reduced if those chestnuts are susceptible to ink disease. Further, the range of this pathogen will extend northward as the climate warms, which may further limit the potential range of the American chestnut. Potassium phosphonate has been found to induce resistance to infection of C. sativa by both inhibiting Phytophthora species directly and by improving the host response, inducing resistance to lesions in phloem tissue and formation of callus. Whether or not this effect would occur in C. dentata is uncertain. Chestnut brown nut rot Brown nut rot is a destructive plant disease caused by the primary agent Gnomoniopsis castaneae and afflicting Castanea species. This pathogen also causes mild disease or exists as an endophyte in other hardwoods. The disease is found in Europe and Oceania. More recently, it has been reported in North America. Chinese gall wasp The Chinese gall wasp attacks all Chestnut species and causes heavy damage. As this species of wasp is a threat to saplings, and is now widely-present in Eastern North American forests, it is a potential problem for reintroduction of the American chestnut. The Chinese parasitoid chalcid wasp Torymus sinensis is considered an effective control method for the Chinese gall wasp. There are now established populations of Torymus sinensis in North America. Reduced population American chestnuts were a common part of the forest canopy in southeast Michigan. Although large trees are currently rare east of the Mississippi River, they exist in pockets in the blight-free West, where the habitat was agreeable for planting: settlers took seeds of American chestnut with them in the 19th century. Huge planted chestnut trees can be found in Sherwood, Oregon, as the Mediterranean climate of the West Coast discourages the fungus, which relies on hot, humid summer weather. American chestnut also thrives as far north as Revelstoke, British Columbia. At present, it is believed that survival of C. dentata for more than a decade in its native range is almost impossible. The fungus uses various oak trees as a host, and while the oak itself is unaffected, American chestnuts nearby will succumb to the blight in approximately a year or more. In addition, the hundreds of chestnut stumps and "living stools" dotting eastern woodlands may still contain active pathogens. It is considered extirpated from Florida and Illinois. The reduced population of American chestnuts directly impacted many species of insects that relied upon the tree species for survival. Of approximately 60 species that feed upon the American chestnut, seven rely entirely on the American chestnut as a food source. Some of these, like the American chestnut moth, are now extinct or severely reduced in population. Attempts at restoration Transgenic blight-resistant American chestnut Researchers at the State University of New York College of Environmental Science and Forestry (SUNY ESF) have developed the Darling 58 chestnut cultivar. This cultivar expresses the gene for wheat oxalate oxidase enzyme, which breaks down the oxalic acid produced by the blight fungus, preventing the death of the tree. When pollen of transgenic fathers fertilizes an ovule of a native mother in the field, those resulting seedlings that express the oxalate oxidase enzyme show growth similar to non-transgenic full siblings, indicating that the transgene does not impede growth under field conditions. The modified chestnut does not affect survival, pollen use, or reproduction of bumble bees. A deregulation petition for the Darling 58 variant has been submitted. If approved, these trees could be the first genetically modified forest trees released in the wild in the United States. American chestnut (C. dentata) is susceptible to ink disease, particularly in the southern part of its natural range. Unlike American chestnut, Japanese chestnut (C. crenata) exhibits resistance to Phytophthora cinnamomi, the fungal pathogen that causes ink disease. The mechanism of resistance of C. crenata to Phytophthora cinnamomi may derive from its expression of the Cast_Gnk2-like gene (99.6% identical with ). Transgenic modification of C. dentata with the Cast_Gnk2-like gene may provide a mechanism for developing American chestnut trees resistant to Phytophthora cinnamomi. Stacking of the Cast_Gnk2-like gene and the oxalate oxidase gene may provide a means of developing genetically modified American chestnut trees resistant to both the chestnut blight and to ink disease. The American Chestnut Research and Restoration Project at SUNY-ESF is not pursuing intellectual property (IP) protection through patents. Dr William Powell, the project's co-director, states that the decision to not pursue a patent on the project's transgenic lines was to allow the plant to be more accessible for conservationists and members of the public. Powell posits that a patent would constrain the spread of the oxalate oxidase transgene into American chestnut populations by limiting the ability to freely plant transgenic trees and cross the trees with surviving American chestnuts or the hybrids produced in the backcross program. Powell states that patents would be a barrier to chestnut restoration and in direct opposition to the program's goals of collaboration. While patent protection is not sought, the non-profit American Chestnut Foundation (TACF) maintains control through a Germplasm Agreement, ensuring authorized use of chestnut germplasm. The agreement safeguards TACF's rights and aligns with the organization's restoration goals. A laboratory error resulted in mistaken use of the Darling 54 cultivar instead of the Darling 58 cultivar in some field trials. TACF and colleagues have also reported decreased growth rates, and poor heritability of resistance of the Darling 58 cultivar. In response, the American Chestnut Foundation withdrew its support of development of the Darling 58 cultivar in December 2023. The American Chestnut Research & Restoration Program at SUNY ESF is continuing to pursue deregulation. In 2022, the SUNY ESF group developed transgenic American chestnut trees incorporating both the oxalate oxidase transgene from wheat and the win3.12 promoter transgene from the necklace poplar. Unlike the CAMV 35S promoter which acts at all times, this poplar promoter drives OxO expression at a low level under basal conditions, but elevates to high levels under conditions of wounding or tissue infection. In laboratory bioassays, win3.12-OxO lines showed elevated disease tolerance similar to that exhibited by blight-resistant Chinese chestnut. Limiting expression of the OxO gene to blight infected tissues is expected to reduce the metabolic cost of gene expression, resulting in increased competitiveness of these new cultivars. As of January 2025, the win3.12 OxO chestnut is still in the research and development phase and has not been deployed for restoration purposes. Intercrossing surviving American chestnuts American Chestnut Cooperators' Foundation (ACCF) is not using crosses with Asian species for blight resistance, but intercrossing among American chestnuts selected for native resistance to the blight, a breeding strategy described by the ACCF as "All-American intercrosses". John Rush Elkins, a research chemist and professor emeritus of chemistry at Concord University, and Gary Griffin, professor of plant pathology at Virginia Tech, think there may be several different characteristics which favor blight resistance. Both Elkins and Griffin have written extensively about the American chestnut. They believe that by making intercrosses among resistant American chestnuts from many locations, they will continue to improve upon the levels of blight resistance to make an American chestnut that can compete in the forest. Griffin, who has been involved with American chestnut restoration for many years, developed a scale for assessing levels of blight resistance, which made it possible to make selections scientifically. He inoculated five-year-old chestnuts with a standard lethal strain of the blight fungus and measured growth of the cankers. Chestnuts with no resistance to blight make rapid-growing, sunken cankers that are deep and kill tissue right to the wood. Resistant chestnuts make slow-growing, swollen cankers that are superficial: live tissue can be recovered under these cankers. The level of blight resistance is judged by periodic measurement of cankers. Grafts from large survivors of the blight epidemic were evaluated following inoculations, and controlled crosses among resistant American chestnut trees were made beginning in 1980. The first "All-American intercrosses" were planted in Virginia Tech's Martin American Chestnut Planting in Giles County, Virginia, and in Beckley, West Virginia. They were inoculated in 1990 and evaluated in 1991 and 1992. Nine of the trees showed resistance equal to their parents, and four of these had resistance comparable to hybrids in the same test. Many ACCF chestnuts have expressed blight resistance equal to or greater than an original blight survivor but so far, only a handful have demonstrated superior, durable blight control. Time will tell if the progeny of these best chestnuts exhibit durable blight resistance in different stress environments. Backcrossing Backcrossing as a treatment for blight was first proposed by Charles Burnham of the University of Minnesota in the 1970s. Burnham, a professor emeritus in agronomy and plant genetics who was considered one of the pioneers of maize genetics, realized that experiments conducted by the USDA to cross-breed American chestnuts with European and Asian chestnuts erroneously assumed that a large number of genes were responsible for blight resistance, while it is currently believed the number of responsible genes is low. The USDA abandoned their cross-breeding program and destroyed local plantings around 1960 after failing to produce a blight-resistant hybrid. Burnham's recognition of the USDA's error led to him joining with others to create The American Chestnut Foundation in 1983, with the sole purpose of breeding a blight-resistant American chestnut. The American Chestnut Foundation is backcrossing blight-resistant Chinese chestnut into American chestnut trees, to recover the American growth characteristics and genetic makeup, and then finally intercrossing the advanced backcross generations to eliminate genes for susceptibility to blight. The first backcrossed American chestnut tree, called "Clapper", survived blight for 25 years, and grafts of the tree have been used by The American Chestnut Foundation since 1983. The Pennsylvania chapter of The American Chestnut Foundation, which seeks to restore the American chestnut to the forests of the Mid-Atlantic states, has planted over 22,000 trees. The Surface Mining Control and Reclamation Act of 1977 requires owners of abandoned coal mines to cover at least 80 percent of their land with vegetation. While many companies planted invasive grasses, others began funding research on planting trees, because they can be more cost-effective, and yield better results. Keith Gilland began planting American chestnut trees in old strip mines in 2008 as a student at Miami University, and to date has planted over 5,000 trees. In 2005, a hybrid tree with mostly American genes was planted on the lawn of the White House. A tree planted in 2005 in the tree library outside the USDA building was still very healthy seven years later; it contains 98% American chestnut DNA and 2% Chinese chestnut DNA. This tree contains enough Chinese chestnut DNA that encodes for systemic resistance genes to resist the blight. This is essential for restoring the American chestnut trees into the Northeast. The Northern Nut Growers Association (NNGA) has also been active in pursuing viable hybrids. From 1962 to 1990, Alfred Szego and other members of the NNGA developed hybrids with Chinese varieties which showed limited resistance. Initially the backcrossing method would breed a hybrid from an American chestnut nut and a Chinese chestnut, the hybrid would then be bred with a normal American chestnut, subsequent breeding would involve a hybrid and an American chestnut or two hybrids, which would increase the genetic makeup of the hybrids primarily American chestnut but still retain the blight resistance of the Chinese chestnut. A backcross breeding program has integrated desirable American chestnut traits with traits from the Chinese chestnut, achieving intermediate resistance to Cryphonectria parasitica and Phytophthora cinnamomi in the hybrid genome. The B3F3 strain, a product of backcrossing and intercrossing with selection for blight resistance, is approximately 94% American chestnut and 6% Chinese chestnut and has been planted experimentally in Maryland in an orchard. Hypovirulence Hypovirus is the only genus in the family Hypoviridae. Members of this genus infect fungal pathogens and reduce their ability to cause disease (hypovirulence). In particular, the virus infects Cryphonectria parasitica, the fungus that causes chestnut blight, which has enabled infected trees to recover from the blight. The use of hypovirulence to control blight originated in Europe where the fungal virus spread naturally through populations of European chestnuts. The reduced ability of the fungus to cause disease allowed the European chestnut to regenerate, creating large stands of trees. Hypovirulence has also been found in North America, but has not spread effectively. The "Arner Tree" of Southern Ontario is one of the best examples of naturally occurring hypovirulence. It is a mature American chestnut that has recovered from severe infections of chestnut blight. The cankers have healed over and the tree continues to grow vigorously. Scientists have discovered that the chestnut blight remaining on the tree is hypovirulent, although isolates taken from the tree do not have the fungal viruses found in other isolates. Trees inoculated with isolates taken from the Arner tree have shown moderate canker control. The cankers of hypovirulent American chestnut trees occur on the outermost tissues of the tree but the cankers do not spread into the growth tissues of the American chestnut tree, thereby providing it with a resistance. Surviving specimens About 2,500 chestnut trees are growing on near West Salem, Wisconsin, which is the world's largest remaining stand of American chestnut. These trees are the descendants of those planted by Martin Hicks, an early settler in the area, who planted fewer than a dozen trees in the late 1800s. Planted outside the natural range of chestnut, these trees escaped the initial onslaught of chestnut blight, but in 1987, scientists found blight in the stand. Scientists are working to try to save the trees. Some 1,348 chestnut trees, varying in size from seedlings to nearly mature trees, are growing in a forest in Western Maine. These chestnuts were originally established in 1982 from four seed-bearing trees sourced from wild stock of a northern Michigan relict population. This grove of trees has dispersed over an area up to 370 meters from the parent trees. The trees appear to be free of chestnut blight. Also in Western Maine, The University of Maine Foundation believes it has the tallest chestnut in North America on its property in Lovell, Maine. As per a Dec. 2015 measurement, it was tall, in girth and believed to be around 100 years old. Two of the largest surviving American chestnut trees are in Jackson County, Tennessee. One, the state champion, has a diameter of and a height of , and the other tree is nearly as large. One of them has been pollinated with hybrid pollen by members of The American Chestnut Foundation; the progeny will have mostly American chestnut genes and some will be blight resistant. On May 18, 2006, a biologist with the Georgia Department of Natural Resources spotted a stand of several trees near Warm Springs, Georgia. One of the trees is approximately 20–30 years old and tall and is the southernmost American chestnut tree known to be flowering and producing nuts. Another large tree was found in Talladega National Forest, Alabama, in June 2005. In the summer of 2007, a stand of trees was discovered near the northeastern Ohio town of Braceville. The stand encompasses four large flowering trees, the largest of which is about tall, sited among hundreds of smaller trees that have not begun to flower, located in and around a sandstone quarry. A combination of factors may account for the survival of these relatively large trees, including low levels of blight susceptibility, hypovirulence, and good site conditions. In particular, some stands may have avoided exposure due to being located at a higher altitude than blighted trees in the neighboring area; the fungal spores are not carried to higher altitudes as easily. In March 2008, officials of the Ohio Department of Natural Resources announced a rare adult American chestnut tree had been discovered in a marsh near Lake Erie. The officials admitted they had known about the tree for seven years, but had kept its existence a secret. The exact location of the tree is still being held secret, both because of the risk of infecting the tree and because an eagle has nested in its branches. They described the tree as being tall and having a circumference of . The American Chestnut Foundation was also only recently told about the tree's existence. Members of the Kentucky chapter of the American Chestnut Foundation have been pollinating a tree found on a farm in Adair County, and a specimen found on Henderson Ridge in Elliott County. The Adair County tree is over one hundred years old. In June 2007, a mature American chestnut was discovered in Farmington, New Hampshire. In rural Missaukee County, Michigan, a blight-free grove of American chestnut trees approximately in size with the largest tree measuring in circumference ( diameter) has been located. It is believed to be the result of nuts planted by early settlers in the area. The American Chestnut Council, located in the local town of Cadillac, Michigan, has verified its identity and existence. Efforts have been initiated to protect the property from logging and development. In Lansing, Michigan, Fenner Nature Center is home to a grove of blight-free American chestnuts descended from the aforementioned grove in Missaukee County. American chestnuts have been located on Beaver Island, a large island in northern Lake Michigan. Hundreds of healthy American chestnuts have been found in the proposed Chestnut Ridge Wilderness Area in the Allegheny National Forest in northwestern Pennsylvania. Many of these trees are large, measuring more than in height. These trees will be protected from logging if the wilderness area, proposed by Friends of Allegheny Wilderness, is passed into law. The Montreal Botanical Garden has the American chestnut among its collection of trees and ornamental shrubs. Three of Portland, Oregon's heritage trees are American chestnuts, along with three Spanish (European) chestnuts. At least two American chestnuts live on the side of Skitchewaug Trail in Springfield, Vermont. Around 300 to 500 trees were spotted in the George Washington National Forest near Augusta County, Virginia, in 2014. Over one dozen trees were at least 12 inches in diameter with several measuring nearly 24 inches in diameter. Only one of the larger trees was a seed and pollen producer with numerous pods and pollen strands lying on ground. The site did, however, have a high presence of chestnut blight, although the seed producing tree and several other large ones were relatively blight-free with minimal to no damage. Two trees were planted 1985, in Nova Scotia, at Dalhousie University, Sexton Campus and are thriving. The donated trees were from saplings grown in Europe, away from the blight. They have diameter trunks and are approximately high. A single mature American chestnut can be found on the front lawn of the McPhail house heritage site in Sault Ste Marie, Ontario, planted by former mayor John Alexander McPhail in the 1920s. Well north of the natural range of the chestnut, it has avoided the blight. There is one American chestnut in Pennsylvania in the county of Columbia in the township of South Centre. It is a hardy, nut-producing tree that has been producing for nearly 30 years. A solitary tree exists in the New York County of Orange, within the Town of Wawayanda. This was planted in the early 1990s as part of a local soil and water conservation district program to identify blight/resistant specimens. It has borne fruit since 2005. A lone but "perfect" American Chestnut tree grows on the Oakdale Campus in Coralville, Iowa. The great majority of chestnut trees in the United States are derived from Dunstan chestnuts, developed in Greensboro, North Carolina, in the 1960s. The Canadian Chestnut Council has a plot growing and harvesting chestnut trees at Tim Hortons Children's Foundation Onondaga Farms. The seedlings are grown at a Simcoe, Ont, Canada station. They are then brought in the spring to this test plantation in St. George, Ontario, between Brantford and Cambridge. Multiple chestnut trees are still alive and nut bearing in Wind River Arboretum, Washington State. Uses Food The nuts were once an important economic resource in North America, being sold on the streets of towns and cities, as they sometimes still are during the Christmas season (usually said to be "roasting on an open fire" because their smell is readily identifiable many blocks away). Chestnuts are edible raw or roasted, though typically preferred roasted. One must peel the brown skin to access the yellowish-white edible portion. The nuts were commonly fed on by various types of wildlife and was also in such a high abundance that they were used to feed livestock by farmers, by allowing those livestock to roam freely into the forests that were predominantly filled with American chestnut trees. The American chestnut tree was important to many Native American tribes in North America as it served as a food source, both for them and the wildlife they hunted, and also as a component in traditional medicine. Furniture and other wood products The January 1888 issue of Orchard and Garden mentions the American chestnut as being "superior in quality to any found in Europe". The wood is straight-grained, strong, and easy to saw and split, and it lacks the radial end grain found on most other hardwoods. The tree was particularly valuable commercially since it grew at a faster rate than oaks. Being rich in tannins, the wood was highly resistant to decay and therefore used for a variety of purposes, including furniture, split-rail fences, shingles, home construction, flooring, piers, plywood, paper pulp, and telephone poles. Tannins were also extracted from the bark for tanning leather. Although larger trees are no longer available for milling, much chestnut wood has been reclaimed from historic barns to be refashioned into furniture and other items. "Wormy" chestnut refers to a defective grade of wood that has insect damage, having been sawn from long-dead, blight-killed trees. This "wormy" wood has since become fashionable for its rustic character.	Biology and health sciences	Nuts	Plants
67742	https://en.wikipedia.org/wiki/Percussion%20cap	Percussion cap	The percussion cap, percussion primer, or caplock, introduced in the early 1820s, is a type of single-use percussion ignition device for muzzle loader firearm locks enabling them to fire reliably in any weather condition. Its invention gave rise to the caplock mechanism or percussion lock system which used percussion caps struck by the hammer to set off the gunpowder charge in rifles and cap and ball firearms. Any firearm using a caplock mechanism is a percussion gun. Any long gun with a cap-lock mechanism and rifled barrel is a percussion rifle. Cap and ball describes cap-lock firearms discharging a single bore-diameter spherical bullet with each shot. Description The percussion cap is a small cylinder of copper or brass with one closed end. Inside the closed end is a small amount of a shock-sensitive explosive material such as mercuric fulminate (discovered in 1800; it was the only practical detonator used from about the mid-19th century to the early 20th century). The caplock mechanism consists of a hammer and a nipple (sometimes referred to as a cone). The nipple contains a hollow conduit which goes into the rearmost part of the gun barrel, and the percussion cap is placed over the nipple hole. Pulling the trigger releases the hammer, which strikes the percussion cap against the nipple (which serves as an anvil), crushes it and detonates the mercury fulminate inside, which releases sparks that travel through the hollow nipple into the barrel and ignite the main powder charge. Percussion caps have been made in small sizes for pistols and larger sizes for rifles and muskets. Origins Earlier firearms used flintlock mechanisms causing a piece of flint to strike a steel frizzen producing sparks to ignite a pan of priming powder and thereby fire the gun's main powder charge. The flintlock mechanism replaced older ignition systems such as the matchlock and wheellock, but all were prone to misfire in wet weather. The discovery of fulminates was made by Edward Charles Howard (1774–1816) in 1800. The invention that made the percussion cap possible using the recently discovered fulminates was patented by the Reverend Alexander John Forsyth of Belhelvie, Aberdeenshire, Scotland, in 1807. The rudimentary percussion system was invented by Forsyth as a solution to the problem that birds would startle when smoke puffed from the powder pan of his flintlock shotgun, giving them sufficient warning to escape the shot. This early percussion lock system operated in a nearly identical fashion to flintlock firearms and used a fulminating primer made of fulminate of mercury, chlorate of potash, sulphur and charcoal, ignited by concussion. His invention of a fulminate-primed firing mechanism deprived the birds of their early warning system, both by avoiding the initial puff of smoke from the flintlock powder pan, as well as shortening the interval between the trigger pull and the shot leaving the muzzle. Forsyth patented his "scent bottle" ignition system in 1807. However, it was not until after Forsyth's patents expired that the conventional percussion cap system was developed. Joseph Manton invented a precursor to the percussion cap in 1814, comprising a copper tube that detonated when crushed. This was further developed in 1822 by the English-born American artist Joshua Shaw, as a copper cup filled with fulminates. The first purpose-built caplock guns were fowling pieces commissioned by sportsmen in Regency era England. Due to the mechanism's compactness and superior reliability compared to the flintlock, gunsmiths were able to manufacture pistols and long guns with two barrels. Early caplock handguns with two or more barrels and a single lock are known as turn-over or twister pistols, due to the need to manually rotate the second barrel to align with the hammer. With the addition of a third barrel, and a ratchet to mechanically turn the barrels while cocking the hammer, these caplock pistols evolved into the pepper-box revolver during the 1830s. The caplock offered many improvements over the flintlock. The caplock was easier and quicker to load, more resilient to weather conditions, and far more reliable than the flintlock. Many of the older flintlock weapons were later converted to the caplock, so that they could take advantage of these features. Parallel developments Joshua Shaw is sometimes credited (primarily by himself) with the development of the first metallic percussion cap in 1814, a reusable one made of iron, then a disposable pewter one in 1815 and finally a copper one in 1816. There is no independent proof of this since Shaw was advised he could not patent it due to Alexander Forsyth's patent for using fulminates to ignite guns being in force between 1807 and 1821. Shaw says he only shared the development of his innovation with a few associates (gunmakers and others) who were sworn to secrecy and never provided affidavits at a later date. Shaw's claim to have been the inventor remains clouded in controversy as he did not patent the idea until 1822, having moved to America in 1817. According to Lewis Winant, the US government's decision to award Shaw $25,000 as compensation for his invention being used by the Army was a mistake. Congress believed Shaw's patent was the earliest in the world and awarded him a large sum of money based on this belief. The investigators had overlooked two French patents and the earlier use of the idea in Britain. The earliest known patent anywhere in the world which specifically mentions a percussion cap and nipple was granted in France on 29 July 1818 to François Prélat, four years before Shaw's patent. Prelat made a habit of copying English patents and inventions and the mode of operation he describes is flawed. Secondly a French patent of a percussion cap and nipple had been granted in 1820 to Deboubert. However predating both of these French claims, the most likely inventor of the percussion cap, according to historian Sidney James Gooding, was Joseph Egg (nephew of Durs Egg), around 1817, . There were other earlier claims. Col. Peter Hawker in 1830 simultaneously claimed and denied being the inventor. "I do not wish to say I was the inventor of it - very probably not" but then immediately recounts that he came up with the idea of simplifying a Manton patch-lock, which could be troublesome, by designing a cap and nipple arrangement around 1816 when the patch lock was patented. He says he then presented a drawing to a reluctant Joseph Manton to make a few copper cap guns which were then sold. Hawker, seems to give Joseph Manton more of the glory eight years later in the 1838 edition of his 'Instructions to young Sportsmen', by stating categorically that "copper tubes and primers were decidedly invented by Joe Manton". By the 1850s Hawker was again claiming the invention for himself in his press advertisements. Despite many years of research by Winant, Gooding and De Witt Bailey, the jury is still out as the competing claims are based on personal accounts and have little or no independently verifiable evidence. While the metal percussion cap was the most popular and widely used type of primer, their small size made them difficult to handle under the stress of combat or while riding a horse. Accordingly, several manufacturers developed alternative, "auto-priming" systems. The "Maynard tape primer", for example, used a roll of paper "caps" much like today's toy cap gun. The Maynard tape primer was fitted to some firearms used in the mid-nineteenth century and a few saw brief use in the American Civil War. Other disc or pellet-type primers held a supply of tiny fulminate detonator discs in a small magazine. Cocking the hammer automatically advanced a disc into position. However, these automatic feed systems were difficult to make with the manufacturing systems in the early and mid-nineteenth century and generated more problems than they solved. They were quickly shelved in favor of a single percussion cap that, while unwieldy in some conditions, could be carried in sufficient quantities to make up for occasionally dropping one, while a jammed tape primer system would instead reduce the rifle to an awkward club. Military firearms This invention was gradually improved, and came to be used, first in a steel cap and then in a copper cap, by various gunmakers and private individuals before coming into general military use nearly thirty years later. The alteration of the military flintlock to the percussion musket was easily accomplished by replacing the powder pan and steel frizzen with a nipple and by replacing the cock or hammer that held the flint by a smaller hammer formed with a hollow made to fit around the nipple when released by the trigger. On the nipple was placed the copper cap containing Shaw's detonating composition of three parts of chlorate of potash, two of fulminate of mercury and one of powdered glass. The hollow in the hammer contained the fragments of the cap if it fragmented, reducing the risk of injury to the firer's eyes. From the 1820s onwards, the armies of Britain, France, Russia, and America began converting their muskets to the new percussion system. Caplocks were generally applied to the British military musket (the Brown Bess) in 1842, a quarter of a century after the invention of percussion powder and after an elaborate government test at Woolwich in 1834. The first percussion firearm produced for the US military was the percussion carbine version (c.1833) of the M1819 Hall rifle. The Americans' breech loading caplock Hall rifles, muzzle loading rifled muskets and Colt Dragoon revolvers gave them an advantage over the smoothbore flintlock Brown Bess muskets used by Santa Anna's troops during the Mexican War. In Japan, matchlock pistols and muskets were converted to percussion from the 1850s onwards, and new guns based on existing designs were manufactured as caplocks. The Austrians instead used a variant of Manton's tube lock in their Augustin musket until the conventional caplock Lorenz rifle was introduced in 1855. The first practical solution for the problem of handling percussion caps in battle was the Prussian 1841 (Dreyse needle gun), which used a long needle to penetrate a paper cartridge filled with black powder and strike the percussion cap that was fastened to the base of the bullet. While it had a number of problems, it was widely used by the Prussians and other German states in the mid-nineteenth century and was a major factor in the 1866 Austro-Prussian War. The needle gun originally fired paper cartridges containing a bullet, powder charge and percussion cap, but by the time of the Franco-Prussian War this had evolved into modern brass ammunition. Later firearms evolution The percussion cap brought about the invention of the modern cartridge case and made possible the general adoption of the breech-loading principle for all varieties of rifles, shotguns and pistols. After the American Civil War, Britain, France, and America began converting existing caplock guns to accept brass rimfire and centrefire cartridges. For muskets such as the 1853 Enfield and 1861 Springfield, this involved installing a firing pin in place of the nipple, and a trapdoor in the breech to accept the new bullets. Examples include the Trapdoor Springfield, Tabatière rifle, Westley Richards and Snider–Enfield conversions. The British Army used Snider Enfields contemporaneously with the Martini–Henry rifle until the .303 bolt action Lee–Metford repeating rifle was introduced in the 1880s. Later, military surplus Sniders were purchased as hunting and defensive weapons by British colonists and trusted local natives. Caplock revolvers such as the Colt Navy and Remington were also widely converted during the late 19th century, by replacing the existing cylinder with one designed for modern ammunition. These were used extensively by the Turks in the Russo-Turkish War, the US Cavalry during the Indian Wars, and also by gunfighters, lawmen, and outlaws in the old west. In the 1840s and 1850s, the percussion cap was first integrated into a metallic cartridge, where the bullet is held in by the casing, the casing is filled with gunpowder, and a primer is placed on the end. By the 1860s and 1870s, breech-loading metallic cartridges had made the percussion cap system obsolete. Today, reproduction percussion firearms are popular for recreational shooters and percussion caps are still available (though some modern muzzleloaders use shotshell primers instead of caps). Most percussion caps now use non-corrosive compounds such as lead styphnate. Other uses Caps are used in cartridges, grenades, rocket-propelled grenades and rescue flares. Percussion caps are also used in land mine fuzes, booby-trap firing devices and anti-handling devices. Most purpose-made military booby-trap firing devices contain some form of spring-loaded firing pin designed to strike a percussion cap connected to a detonator at one end. The detonator is inserted into an explosive charge—e.g., C-4 or a block of TNT. Triggering the booby-trap (e.g., by pulling on a trip-wire) releases the cocked firing pin that flips forward to strike the percussion cap, firing it and the attached detonator; the shock-wave from the detonator sets off the main explosive charge.	Technology	Mechanisms_2	null
67782	https://en.wikipedia.org/wiki/Sailing%20ship	Sailing ship	A sailing ship is a sea-going vessel that uses sails mounted on masts to harness the power of wind and propel the vessel. There is a variety of sail plans that propel sailing ships, employing square-rigged or fore-and-aft sails. Some ships carry square sails on each mast—the brig and full-rigged ship, said to be "ship-rigged" when there are three or more masts. Others carry only fore-and-aft sails on each mast, for instance some schooners. Still others employ a combination of square and fore-and-aft sails, including the barque, barquentine, and brigantine. Early sailing ships were used for river and coastal waters in Ancient Egypt and the Mediterranean. The Austronesian peoples developed maritime technologies that included the fore-and-aft crab-claw sail and with catamaran and outrigger hull configurations, which enabled the Austronesian expansion into the islands of the Indo-Pacific. This expansion originated in Taiwan BC and propagated through Island Southeast Asia, reaching Near Oceania BC, Hawaii AD, and New Zealand AD. The maritime trading network in the Indo-Pacific dates from at least 1500 BC. Later developments in Asia produced the junk and dhow—vessels that incorporated features unknown in Europe at the time. European sailing ships with predominantly square rigs became prevalent during the Age of Discovery (15th to 17th centuries), when they crossed oceans between continents and around the world. In the European Age of Sail, a full-rigged ship was one with a bowsprit and three masts, each of which consists of a lower, top, and topgallant mast. Most sailing ships were merchantmen, but the Age of Sail also saw the development of large fleets of well-armed warships. The many steps of technological development of steamships during the 19th century provided slowly increasing competition for sailing ships — initially only on short routes where high prices could be charged. By the 1880s, ships with triple-expansion steam engines had the fuel efficiency to compete with sail on all major routes — and with scheduled sailings that were not affected by the wind direction. However, commercial sailing vessels could still be found working into the 20th century, although in reducing numbers and only in certain trades. History By the time of the Age of Discovery—starting in the 15th century—square-rigged, multi-masted vessels were the norm and were guided by navigation techniques that included the magnetic compass and making sightings of the sun and stars that allowed transoceanic voyages. The Age of Sail reached its peak in the 18th and 19th centuries with large, heavily armed battleships and merchant sailing ships. Sailing and steam ships coexisted for much of the 19th century. The steamers of the early part of the century had very poor fuel efficiency and were suitable only for a small number of roles, such as towing sailing ships and providing short route passenger and mail services. Both sailing and steam ships saw large technological improvements over the century. Ultimately the two large stepwise improvements in fuel efficiency of compound and then triple-expansion steam engines made the steamship, by the 1880s, able to compete in the vast majority of trades. Commercial sail still continued into the 20th century, with the last ceasing to trade by . South China Sea and Austronesia Early sea-going sailing vessels were used by the Austronesian peoples. The invention of catamarans, outriggers, and crab claw sails enabled the Austronesian Expansion at around 3000 to 1500 BC. From Taiwan, they rapidly colonized the islands of Maritime Southeast Asia, then sailed further onwards to Micronesia, Island Melanesia, Polynesia, and Madagascar. Austronesian rigs were distinctive in that they had spars supporting both the upper and lower edges of the sails (and sometimes in between), in contrast to western rigs which only had a spar on the upper edge. Large Austronesian trading ships with as many as four sails were recorded by Han dynasty (206 BC – 220 AD) scholars as the kunlun bo or K'un-lun po (崑崙舶, lit. "ship of the Kunlun people"). They were booked by Chinese Buddhist pilgrims for passage to Southern India and Sri Lanka. Bas reliefs of large Javanese outriggers ships with various configurations of tanja sails are also found in the Borobudur temple, dating back to the 8th century CE. By the 10th century AD, the Song dynasty started building the first Chinese seafaring junks, which adopted several features of the K'un-lun po. The junk rig in particular, became associated with Chinese coast-hugging trading ships. Junks in China were constructed from teak with pegs and nails; they featured watertight compartments and acquired center-mounted tillers and rudders. These ships became the basis for the development of Chinese warships during the Mongol Yuan dynasty, and were used in the unsuccessful Mongol invasions of Japan and Java. The Ming dynasty (1368–1644) saw the use of junks as long-distance trading vessels. Chinese Admiral Zheng He reportedly sailed to India, Arabia, and southern Africa on a trade and diplomatic mission. Literary lore suggests that his largest vessel, the "Treasure Ship", measured in length and in width, whereas modern research suggests that it was unlikely to have exceeded in length. Mediterranean and Baltic Sailing ships in the Mediterranean region date back to at least 3000 BC, when Egyptians used a bipod mast to support a single square sail on a vessel that mainly relied on multiple paddlers. Later the mast became a single pole, and paddles were supplanted with oars. Such vessels plied both the Nile and the Mediterranean coast. The Minoan civilization of Crete may have been the world's first thalassocracy brought to prominence by sailing vessels dating to before 1800 BC (Middle Minoan IIB). Between 1000 BC and 400 AD, the Phoenicians, Greeks and Romans developed ships that were powered by square sails, sometimes with oars to supplement their capabilities. Such vessels used a steering oar as a rudder to control direction. Starting in the 8th century in Denmark, Vikings were building clinker-constructed longships propelled by a single, square sail, when practical, and oars, when necessary. A related craft was the knarr, which plied the Baltic and North Seas, using primarily sail power. The windward edge of the sail was stiffened with a beitass, a pole that fitted into the lower corner of the sail, when sailing close to the wind. Indian Ocean India's maritime history began during the 3rd millennium BCE when inhabitants of the Indus Valley initiated maritime trading contact with Mesopotamia. Indian kingdoms such as the Kalinga from as early as 2nd century CE are believed to have had sailing ships. One of the earliest instances of documented evidence of Indian sailing ship building comes from the mural of three-masted ship in the Ajanta caves that date back to 400-500 CE. The Indian Ocean was the venue for increasing trade between India and Africa between 1200 and 1500. The vessels employed would be classified as dhows with lateen rigs. During this interval such vessels grew in capacity from 100 to 400 tonnes. Dhows were often built with teak planks from India and Southeast Asia, sewn together with coconut husk fiber—no nails were employed. This period also saw the implementation of center-mounted rudders, controlled with a tiller. Global exploration Technological advancements that were important to the Age of Discovery in the 15th century were the adoption of the magnetic compass and advances in ship design. The compass was an addition to the ancient method of navigation based on sightings of the sun and stars. The compass was invented by Chinese. It had been used for navigation in China by the 11th century and was adopted by the Arab traders in the Indian Ocean. The compass spread to Europe by the late 12th or early 13th century. Use of the compass for navigation in the Indian Ocean was first mentioned in 1232. The Europeans used a "dry" compass, with a needle on a pivot. The compass card was also a European invention. At the beginning of the 15th century, the carrack was the most capable European ocean-going ship. It was carvel-built and large enough to be stable in heavy seas. It was capable of carrying a large cargo and the provisions needed for very long voyages. Later carracks were square-rigged on the foremast and mainmast and lateen-rigged on the mizzenmast. They had a high rounded stern with large aftcastle, forecastle and bowsprit at the stem. As the predecessor of the galleon, the carrack was one of the most influential ship designs in history; while ships became more specialized in the following centuries, the basic design remained unchanged throughout this period. Ships of this era were only able to sail approximately 70° into the wind and tacked from one side to the other across the wind with difficulty, which made it challenging to avoid shipwrecks when near shores or shoals during storms. Nonetheless, such vessels reached India around Africa with Vasco da Gama, the Americas with Christopher Columbus, and around the world under Ferdinand Magellan. 1700 to 1850 Sailing ships became longer and faster over time, with ship-rigged vessels carrying taller masts with more square sails. Other sail plans emerged, as well, that had just fore-and-aft sails (schooners), or a mixture of the two (brigantines, barques and barquentines). Warships Cannons were introduced in the 14th century, but did not become common at sea until they could be reloaded quickly enough to be reused in the same battle. The size of a ship required to carry a large number of cannon made oar-based propulsion impossible, and warships came to rely primarily on sails. The sailing man-of-war emerged during the 16th century. By the middle of the 17th century, warships were carrying increasing numbers of cannon on three decks. Naval tactics evolved to bring each ship's firepower to bear in a line of battle—coordinated movements of a fleet of warships to engage a line of ships in the enemy fleet. Carracks with a single cannon deck evolved into galleons with as many as two full cannon decks, which evolved into the man-of-war, and further into the ship of the line—designed for engaging the enemy in a line of battle. One side of a ship was expected to shoot broadsides against an enemy ship at close range. In the 18th century, the small and fast frigate and sloop-of-war—too small to stand in the line of battle—evolved to convoy trade, scout for enemy ships and blockade enemy coasts. Clippers The term "clipper" started to be used in the first quarter of the 19th century. It was applied to sailing vessels designed primarily for speed. Only a small proportion of sailing vessels could properly have the term applied to them. Early examples were the schooners and brigantines, called Baltimore clippers, used for blockade running or as privateers in the War of 1812 and afterwards for smuggling opium or illegally transporting slaves. Larger clippers, usually ship or barque rigged and with a different hull design, were built for the California trade (from east coast USA ports to San Francisco) after gold was discovered in 1848 the associated shipbuilding boom lasted until 1854. Clippers were built for trade between the United Kingdom and China after the East India Company lost its monopoly in 1834. The primary cargo was tea, and sailing ships, particularly tea clippers, dominated this long-distance route until the development of fuel efficient steamships coincided with the opening of the Suez Canal in 1869. Other clippers worked on the Australian immigrant routes or, in smaller quantities, in any role where a fast passage secured higher rates of freight or passenger fares. Whilst many clippers were ship rigged, the definition is not limited to any rig. Clippers were generally built for a specific trade: those in the California trade had to withstand the seas of Cape Horn, whilst Tea Clippers were designed for the lighter and contrary winds of the China Sea. All had fine lines, with a well streamlined hull and carried a large sail area. To get the best of this, a skilled and determined master was needed in command. Copper sheathing During the Age of Sail, ships' hulls were under frequent attack by shipworm (which affected the structural strength of timbers), and barnacles and various marine weeds (which affected ship speed). Since before the common era, a variety of coatings had been applied to hulls to counter this effect, including pitch, wax, tar, oil, sulfur and arsenic. In the mid 18th century copper sheathing was developed as a defense against such bottom fouling. After coping with problems of galvanic deterioration of metal hull fasteners, sacrificial anodes were developed, which were designed to corrode, instead of the hull fasteners. The practice became widespread on naval vessels, starting in the late 18th century, and on merchant vessels, starting in the early 19th century, until the advent of iron and steel hulls. 1850 to 1900 Iron-hulled sailing ships, often referred to as "windjammers" or "tall ships", represented the final evolution of sailing ships at the end of the Age of Sail. They were built to carry bulk cargo for long distances in the nineteenth and early twentieth centuries. They were the largest of merchant sailing ships, with three to five masts and square sails, as well as other sail plans. They carried lumber, guano, grain or ore between continents. Later examples had steel hulls. Iron-hulled sailing ships were mainly built from the 1870s to 1900, when steamships began to outpace them economically, due to their ability to keep a schedule regardless of the wind. Steel hulls also replaced iron hulls at around the same time. Even into the twentieth century, sailing ships could hold their own on transoceanic voyages such as Australia to Europe, since they did not require bunkerage for coal nor fresh water for steam, and they were faster than the early steamers, which usually could barely make . The four-masted, iron-hulled ship, introduced in 1875 with the full-rigged , represented an especially efficient configuration that prolonged the competitiveness of sail against steam in the later part of the 19th century. The largest example of such ships was the five-masted, full-rigged ship , which had a load capacity of 7,800 tonnes. Ships transitioned from all sail to all steam-power from the mid 19th century into the 20th. Five-masted Preussen used steam power for driving the winches, hoists and pumps, and could be manned by a crew of 48, compared with four-masted Kruzenshtern, which has a crew of 257. Coastal top-sail schooners with a crew as small as two managing the sail handling became an efficient way to carry bulk cargo, since only the fore-sails required tending while tacking and steam-driven machinery was often available for raising the sails and the anchor. 1950 to 2000 In the 20th century, the DynaRig allowed central, automated control of all sails in a manner that obviates the need for sending crew aloft. This was developed in the 1960s in Germany as a low-carbon footprint propulsion alternative for commercial ships. The rig automatically sets and reefs sails; its mast rotates to align the sails with the wind. The sailing yachts Maltese Falcon and Black Pearl employ the rig. 21st century and contemporary experimental sail In the 21st century, due to concern about climate change and the possibility of cost savings, companies explored using wind-power to reduce heavy fuel needs on large containerized cargo ships. By 2023, around 30 ships were using sails or attached kites, with the number expected to grow. The following year, The Economist wrote that the technology was at an inflection point as it moved from trials and testing towards adoption by the industry. Features Every sailing ship has a sail plan that is adapted to the purpose of the vessel and the ability of the crew; each has a hull, rigging and masts to hold up the sails that use the wind to power the ship; the masts are supported by standing rigging and the sails are adjusted by running rigging. Hull Hull shapes for sailing ships evolved from being relatively short and blunt to being longer and finer at the bow. By the nineteenth century, ships were built with reference to a half model, made from wooden layers that were pinned together. Each layer could be scaled to the actual size of the vessel in order to lay out its hull structure, starting with the keel and leading to the ship's ribs. The ribs were pieced together from curved elements, called futtocks and tied in place until the installation of the planking. Typically, planking was caulked with a tar-impregnated yarn made from manila or hemp to make the planking watertight. Starting in the mid-19th century, iron was used first for the hull structure and later for its watertight sheathing. Masts Until the mid-19th century all vessels' masts were made of wood formed from a single or several pieces of timber which typically consisted of the trunk of a conifer tree. From the 16th century, vessels were often built of a size requiring masts taller and thicker than could be made from single tree trunks. On these larger vessels, to achieve the required height, the masts were built from up to four sections (also called masts), known in order of rising height above the decks as the lower, top, topgallant and royal masts. Giving the lower sections sufficient thickness necessitated building them up from separate pieces of wood. Such a section was known as a made mast, as opposed to sections formed from single pieces of timber, which were known as pole masts. Starting in the second half of the 19th century, masts were made of iron or steel. For ships with square sails the principal masts, given their standard names in bow to stern (front to back) order, are: Fore-mast — the mast nearest the bow, or the mast forward of the main-mast with sections: fore-mast lower, fore topmast, and fore topgallant mast Main-mast — the tallest mast, usually located near the center of the ship with sections: main-mast lower, main topmast, main topgallant mast, royal mast (sometimes) Mizzen-mast — the aft-most mast. Typically shorter than the fore-mast with sections: mizzen-mast lower, mizzen topmast, and mizzen topgallant mast. Sails Each rig is configured in a sail plan, appropriate to the size of the sailing craft. Both square-rigged and fore-and-aft rigged vessels have been built with a wide range of configurations for single and multiple masts. Types of sail that can be part of a sail plan can be broadly classed by how they are attached to the sailing craft: To a stay — Sails attached to stays, include jibs, which are attached to forestays and staysails, which are mounted on other stays (typically wire cable) that support other masts from the bow aft. To a mast — Fore-and-aft sails directly attached to the mast at the luff include gaff-rigged quadrilateral and Bermuda triangular sails. To a spar — Sails attached to a spar include both square sails and such fore-and-aft quadrilateral sails as lug rigs, junk and spritsails and such triangular sails as the lateen, and the crab claw. Rigging Sailing ships have standing rigging to support the masts and running rigging to raise the sails and control their ability to draw power from the wind. The running rigging has three main roles, to support the sail structure, to shape the sail and to adjust its angle to the wind. Square-rigged vessels require more controlling lines than fore-and-aft rigged ones. Standing rigging Sailing ships prior to the mid-19th century used wood masts with hemp-fiber standing rigging. As rigs became taller by the end of the 19th century, masts relied more heavily on successive spars, stepped one atop the other to form the whole, from bottom to top: the lower mast, top mast, and topgallant mast. This construction relied heavily on support by a complex array of stays and shrouds. Each stay in either the fore-and-aft or athwartships direction had a corresponding one in the opposite direction providing counter-tension. Fore-and-aft the system of tensioning started with the stays that were anchored in front each mast. Shrouds were tensioned by pairs of deadeyes, circular blocks that had the large-diameter line run around them, whilst multiple holes allowed smaller line—lanyard—to pass multiple times between the two and thereby allow tensioning of the shroud. After the mid-19th century square-rigged vessels were equipped with iron wire standing rigging, which was superseded with steel wire in the late 19th century. Running rigging Halyards, used to raise and lower the yards, are the primary supporting lines. In addition, square rigs have lines that lift the sail or the yard from which it is suspended that include: brails, buntlines, lifts and leechlines. Bowlines and clew lines shape a square sail. To adjust the angle of the sail to wind braces are used to adjust the fore and aft angle of a yard of a square sail, while sheets attach to the clews (bottom corners) of a sail to control the sail's angle to the wind. Sheets run aft, whereas tacks are used to haul the clew of a square sail forward. Crew The crew of a sailing ship is divided between officers (the captain and his subordinates) and seamen or ordinary hands. An able seaman was expected to "hand, reef, and steer" (handle the lines and other equipment, reef the sails, and steer the vessel). The crew is organized to stand watch—the oversight of the ship for a period—typically four hours each. Richard Henry Dana Jr. and Herman Melville each had personal experience aboard sailing vessels of the 19th century. Merchant vessel Dana described the crew of the merchant brig, Pilgrim, as comprising six to eight common sailors, four specialist crew members (the steward, cook, carpenter and sailmaker), and three officers: the captain, the first mate and the second mate. He contrasted the American crew complement with that of other nations on whose similarly sized ships the crew might number as many as 30. Larger merchant vessels had larger crews. Warship Melville described the crew complement of the frigate warship, United States, as about 500—including officers, enlisted personnel and 50 Marines. The crew was divided into the starboard and larboard watches. It was also divided into three tops, bands of crew responsible for setting sails on the three masts; a band of sheet-anchor men, whose station was forward and whose job was to tend the fore-yard, anchors and forward sails; the after guard, who were stationed aft and tended the mainsail, spanker and manned the various sheets, controlling the position of the sails; the waisters, who were stationed midships and had menial duties attending the livestock, etc.; and the holders, who occupied the lower decks of the vessel and were responsible for the inner workings of the ship. He additionally named such positions as, boatswains, gunners, carpenters, coopers, painters, tinkers, stewards, cooks and various boys as functions on the man-of-war. 18-19th century ships of the line had a complement as high as 850. Ship handling Handling a sailing ship requires management of its sails to power—but not overpower—the ship and navigation to guide the ship, both at sea and in and out of harbors. Under sail Key elements of sailing a ship are setting the right amount of sail to generate maximum power without endangering the ship, adjusting the sails to the wind direction on the course sailed, and changing tack to bring the wind from one side of the vessel to the other. Setting sail A sailing ship crew manages the running rigging of each square sail. Each sail has two sheets that control its lower corners, two braces that control the angle of the yard, two clewlines, four buntlines and two reef tackles. All these lines must be manned as the sail is deployed and the yard raised. They use a halyard to raise each yard and its sail; then they pull or ease the braces to set the angle of the yard across the vessel; they pull on sheets to haul lower corners of the sail, clews, out to yard below. Under way, the crew manages reef tackles, haul leeches, reef points, to manage the size and angle of the sail; bowlines pull the leading edge of the sail (leech) taut when close hauled. When furling the sail, the crew uses clewlines, haul up the clews and buntlines to haul up the middle of sail up; when lowered, lifts support each yard. In strong winds, the crew is directed to reduce the number of sails or, alternatively, the amount of each given sail that is presented to the wind by a process called reefing. To pull the sail up, seamen on the yardarm pull on reef tackles, attached to reef cringles, to pull the sail up and secure it with lines, called reef points. Dana spoke of the hardships of sail handling during high wind and rain or with ice covering the ship and its rigging. Changing tack Sailing vessels cannot sail directly into the wind. Instead, square-riggers must sail a course that is between 60° and 70° away from the wind direction and fore-and aft vessels can typically sail no closer than 45°. To reach a destination, sailing vessels may have to change course and allow the wind to come from the opposite side in a procedure, called tacking, when the wind comes across the bow during the maneuver. When tacking, a square-rigged vessel's sails must be presented squarely to the wind and thus impede forward motion as they are swung around via the yardarms through the wind as controlled by the vessel's running rigging, using braces—adjusting the fore and aft angle of each yardarm around the mast—and sheets attached to the clews (bottom corners) of each sail to control the sail's angle to the wind. The procedure is to turn the vessel into the wind with the hind-most fore-and-aft sail (the spanker), pulled to windward to help turn the ship through the eye of the wind. Once the ship has come about, all the sails are adjusted to align properly with the new tack. Because square-rigger masts are more strongly braced from behind than from ahead, tacking is a dangerous procedure in strong winds; the ship may lose forward momentum (become caught in stays) and the rigging may fail from the wind coming from ahead. The ship may also lose momentum at wind speeds of less than . Under these conditions, the choice may be to wear ship—to turn the ship away from the wind and around 240° onto the next tack (60° off the wind). A fore-and-aft rig permits the wind to flow past the sail, as the craft head through the eye of the wind. Most rigs pivot around a stay or the mast, while this occurs. For a jib, the old leeward sheet is released as the craft heads through the wind and the old windward sheet is tightened as the new leeward sheet to allow the sail to draw wind. Mainsails are often self-tending and slide on a traveler to the opposite side. On certain rigs, such as lateens and luggers, the sail may be partially lowered to bring it to the opposite side. Navigation Early navigational techniques employed observations of the sun, stars, waves and birdlife. In the 15th century, the Chinese were using the magnetic compass to identify direction of travel. By the 16th century in Europe, navigational instruments included the quadrant, the astrolabe, cross staff, dividers and compass. By the time of the Age of Exploration these tools were being used in combination with a log to measure speed, a lead line to measure soundings, and a lookout to identify potential hazards. Later, an accurate marine sextant became standard for determining latitude and was used with an accurate chronometer to calculate longitude. Passage planning begins with laying out a route along a chart, which comprises a series of courses between fixes—verifiable locations that confirm the actual track of the ship on the ocean. Once a course has been set, the person at the helm attempts to follow its direction with reference to the compass. The navigator notes the time and speed at each fix to estimate the arrival at the next fix, a process called dead reckoning. For coast-wise navigation, sightings from known landmarks or navigational aids may be used to establish fixes, a process called pilotage. At sea, sailing ships used celestial navigation on a daily schedule, as follows: Continuous dead reckoning plot Star observations at morning twilight for a celestial fix Morning Sun observation to determine compass error by azimuth observation of the Sun Noontime observation of the Sun for noon latitude line for determination the day's run and day's set and drift Afternoon sun line to determine compass error by azimuth observation of the Sun Star observations at evening twilight for a celestial fix Fixes were taken with a marine sextant, which measures the distance of the celestial body above the horizon. Entering and leaving harbor Given the limited maneuverability of sailing ships, it could be difficult to enter and leave harbor with the presence of a tide without coordinating arrivals with a flooding tide and departures with an ebbing tide. In harbor, a sailing ship stood at anchor, unless it needed to be loaded or unloaded at a dock or pier, in which case it might be warped alongside or towed by a tug. Warping involved using a long rope (the warp) between the ship and a fixed point on the shore. This was pulled on by a capstan on shore, or on the ship. This might be a multi-stage process if the route was not simple. If no fixed point was available, a kedge anchor might be taken out in a ship's boat to a suitable point and the ship then pulled up to the kedge. Square rigged vessels could use backing and filling (of the sails) to manoeuvre in a tideway, or control could be maintained by drudging the anchor - lower the anchor until it touches the bottom so that the dragging anchor gives steerage way in the flow of the tide. Examples These are examples of sailing ships; some terms have multiple meanings: Defined by general configuration Caravel: small maneuverable ship, lateen rigged Carrack: three or four masted ship, square-rigged forward, lateen-rigged aft Clipper: a merchant ship designed specifically for speed Cog: plank-built, one-masted, square-rigged vessel Dhow: a lateen-rigged merchant or fishing vessel Djong: large tradeship used by ancient Indonesian and Malaysian people Fluyt: a Dutch oceangoing merchant vessel, rigged similarly to a galleon Galleon: a large, primarily square-rigged, armed cargo carrier of the sixteenth and seventeenth centuries Junk: a lug-rigged Chinese ship, which included many types, models and variants. Koch: small, Russian clinker-built ship, designed for use in Arctic waters Longship: vessels used by the Vikings, with a single mast and square sail, also propelled by oars. Pinisi: Indonesia's traditional sailing ship Pink: in the Atlantic, a small oceangoing ship with a narrow stern. Snow: a brig carrying a square mainsail and often a spanker on a trysail mast Sailing superyacht: a large sailing yacht Waʻa kaulua: Polynesian double-hulled voyaging canoe Windjammer: (informal) large merchant sailing ship with an iron or steel hull Defined by sail plan All masts have fore-and-aft sails Schooner: fore-and-aft rigged sails, with two or more masts, the aftermost mast taller or equal to the height of the forward All masts have square sails Brig: two masts, square rigged (may have a spanker on the aftermost) Full-rigged ship: three or more masts, all of them square rigged Mixture of masts with square sails and masts with fore-and-aft sails Barque, or "bark": at least three masts, fore-and-aft rigged mizzen mast Barquentine: at least three masts with all but the foremost fore-and-aft rigged Bilander: a ship or brig with a lug-rigged mizzen sail Brigantine: two masts, with the foremast square-rigged Hermaphrodite brig: a brigantine Military vessels Corvette: lightly armed, fast sailing vessel Cutter: small naval vessel, fore-and-aft rigged, single mast with two headsails Frigate: a ship-rigged warship with a single gundeck Ship of the line: the largest warship in European navies, ship-rigged Xebec: a Mediterranean warship adapted from a galley, with three lateen-rigged masts	Technology	Maritime transport	null
67783	https://en.wikipedia.org/wiki/Sailboat	Sailboat	A sailboat or sailing boat is a boat propelled partly or entirely by sails and is smaller than a sailing ship. Distinctions in what constitutes a sailing boat and ship vary by region and maritime culture. Types Although sailboat terminology has varied across history, many terms have specific meanings in the context of modern yachting. A great number of sailboat-types may be distinguished by size, hull configuration, keel type, purpose, number and configuration of masts, and sail plan. Popular monohull designs include: Cutter The cutter is similar to a sloop with a single mast and mainsail, but generally carries the mast further aft to allow for two foresails, a jib and staysail, to be attached to the head stay and inner forestay, respectively. Once a common racing configuration, today it gives versatility to cruising boats, especially in allowing a small staysail to be flown from the inner stay in high winds. Catboat A catboat has a single mast mounted far forward and does not carry a jib. Most modern designs have only one sail, the mainsail; however, the traditional catboat could carry multiple sails from the gaff rig. Catboat is a charming and distinctive sailboat featuring a single mast with a single large sail, known as a gaff-rigged sail, and a broad beam that ensures stability. This type of vessel, named after the "cat" tackle used in sailing, has a rich history dating back to the 19th century in the coastal regions of the United States, particularly New England, where it was widely used by fishermen and sailors. With its straightforward design and uncomplicated rigging, the catboat offers a straightforward and laid-back sailing experience, making it an ideal choice for beginners and pleasure sailors alike. Even today, catboats continue to be cherished by enthusiasts who appreciate their heritage and enjoy their picturesque appearance while cruising through the waterways. Dinghy A dinghy is a type of small open sailboat commonly used for recreation, sail training, and tending a larger vessel. They are popular in youth sailing programs for their short LOA, simple operation and minimal maintenance. They have three (or fewer) sails: the mainsail, jib, and spinnaker. Ketch Ketches are similar to a sloop, but there is a second shorter mast astern of the mainmast, but forward of the rudder post. The second mast is called the mizzen mast and the sail is called the mizzen sail. A ketch can also be Cutter-rigged with two head sails. Schooner A schooner has a mainmast taller than its foremast, distinguishing it from a ketch or a yawl. A schooner can have more than two masts, with the foremast always lower than the foremost main. Traditional topsail schooners have topmasts allowing triangular topsails sails to be flown above their gaff sails; many modern schooners are Bermuda rigged. Sloop The most common modern sailboat is the sloop, which features one mast and two sails, typically a Bermuda rigged main, and a headsail. This simple configuration is very efficient for sailing into the wind. A fractional rigged sloop has its forestay attached at a point below the top of the mast, allowing the mainsail to be flattened to improve performance by raking the upper part of the mast aft by tensioning the backstay. A smaller headsail is easier for a short-handed crew to manage. Yawl A yawl is similar to a ketch, with a shorter mizzen mast carried astern the rudderpost more for balancing the helm than propulsion. Hulls Traditional sailboats are monohulls, but multi-hull catamarans and trimarans are gaining popularity. Monohull boats generally rely on ballast for stability and usually are displacement hulls. This stabilizing ballast can, in boats designed for racing, be as much as 50% of the weight of the boat, but is generally around 30%. It creates two problems; one, it gives the monohull tremendous inertia, making it less maneuverable and reducing its acceleration. Secondly, unless it has been built with buoyant foam or air tanks, if a monohull fills with water, it will sink. Multihulls rely on the geometry and the broad stance of their multiple hulls for their stability, eschewing any form of ballast. Some multihulls are designed to be as light-weight as possible while still maintaining structural integrity. They can be built with foam-filled flotation chambers and some modern trimarans are rated as unsinkable, meaning that, should every crew compartment be completely filled with water, the hull itself has sufficient buoyancy to remain afloat. A multihull optimized for light weight (at the expense of cruising amenities and storage for food and other supplies), combined with the absence of ballast can result in performance gains in terms of acceleration, top speed, and manoeuvrability. The lack of ballast makes it much easier to get a lightweight multihull on plane, reducing its wetted surface area and thus its drag. Reduced overall weight means a reduced draft, with a much reduced underwater profile. This, in turn, results directly in reduced wetted surface area and drag. Without a ballast keel, multihulls can go in shallow waters where monohulls can not. There are trade-offs, however, in multihull design. A well designed ballasted boat can recover from a capsize, even from turning over completely. Righting a multihull that has gotten upside down is difficult in any case and impossible without help unless the boat is small or carries special equipment for the purpose. Multihulls often prove more difficult to tack, since the reduced weight leads directly to reduced momentum, causing multihulls to more quickly lose speed when headed into the wind. Also, structural integrity is much easier to achieve in a one piece monohull than in a two or three piece multihull whose connecting structure must be substantial and well connected to the hulls. All these hull types may also be manufactured as, or outfitted with, hydrofoils. Keel All vessels have a keel, it is the backbone of the hull. In traditional construction, it is the structure upon which all else depends. Modern monocoque designs include a virtual keel. Even multihulls have keels. On a sailboat, the word "keel" is also used to refer to the area that is added to the hull to improve its lateral plane. The lateral plane is what prevents leeway and allows sailing towards the wind. This can be an external piece or a part of the hull. Most monohulls larger than a dinghy require built-in ballast. Depending on the design of the boat, ballast may be 20 to 50 percent of the displacement. The ballast is often integrated into their keels as large masses of lead or cast iron. This secures the ballast and gets it as low as possible to improve its effectiveness. External keels are cast in the shape of the keel. A monohull's keel is made effective by a combination of weight, depth, and length. Most modern monohull boats have fin keels, which are heavy and deep, but short in relation to the hull length. More traditional yachts carried a full keel which is generally half or more of the length of the boat. A recent feature is a winged keel, which is short and shallow, but carries a lot of weight in two "wings" which run sideways from the main part of the keel. Even more recent is the concept of canting keels, designed to move the weight at the bottom of a sailboat to the upwind side, allowing the boat to carry more sails. A twin keel has the benefit of a shallower draft and can allow the boat to stand on dry land. Multihulls, on the other hand, have minimal need for such ballast, as they depend on the geometry of their design, the wide base of their multiple hulls, for their stability. Designers of performance multihulls, such as the Open 60's, go to great lengths to reduce overall boat weight as much as possible. This leads some to comment that designing a multihull is similar to designing an aircraft. A centreboard or daggerboard is retractable lightweight keel which can be pulled up in shallow water. Mast On small sailboats, masts may be "stepped" (put in place) with the bottom end in a receptacle that is supported above the keel of the boat or on the deck or other superstructure that allows the mast to be raised at a hinge point until it is erect. Some masts are supported solely at the keel and laterally at the deck and are called "unstayed". Most masts rely in part or entirely (for those stepped on the deck) on standing rigging, supporting them side-to-side and fore-and aft to hold them up. Masts over may require a crane and are typically stepped on the keel through any cabin or other superstructure. Auxiliary propulsion Many sailboats have an alternate means of propulsion, in case the wind dies or where close maneuvering under sail is impractical. The smallest boats may use a paddle; bigger ones may have oars; still others may employ an outboard motor, mounted on the transom; still others may have an inboard engine.	Technology	Naval transport	null
67941	https://en.wikipedia.org/wiki/Cassini%E2%80%93Huygens	Cassini–Huygens	Cassini–Huygens ( ), commonly called Cassini, was a space-research mission by NASA, the European Space Agency (ESA), and the Italian Space Agency (ASI) to send a space probe to study the planet Saturn and its system, including its rings and natural satellites. The Flagship-class robotic spacecraft comprised both NASA's Cassini space probe and ESA's Huygens lander, which landed on Saturn's largest moon, Titan. Cassini was the fourth space probe to visit Saturn and the first to enter its orbit, where it stayed from 2004 to 2017. The two craft took their names from the astronomers Giovanni Cassini and Christiaan Huygens. Launched aboard a Titan IVB/Centaur on October 15, 1997, Cassini was active in space for nearly 20 years, with 13 years spent orbiting Saturn and studying the planet and its system after entering orbit on July 1, 2004. The voyage to Saturn included flybys of Venus (April 1998 and July 1999), Earth (August 1999), the asteroid 2685 Masursky, and Jupiter (December 2000). The mission ended on September 15, 2017, when Cassinis trajectory took it into Saturn's upper atmosphere and it burned up in order to prevent any risk of contaminating Saturn's moons, which might have offered habitable environments to stowaway terrestrial microbes on the spacecraft. The mission was successful beyond expectations – NASA's Planetary Science Division Director, Jim Green, described Cassini-Huygens as a "mission of firsts" that has revolutionized human understanding of the Saturn system, including its moons and rings, and our understanding of where life might be found in the Solar System. Cassinis planners originally scheduled a mission of four years, from June 2004 to May 2008. The mission was extended for another two years until September 2010, branded the Cassini Equinox Mission. The mission was extended a second and final time with the Cassini Solstice Mission, lasting another seven years until September 15, 2017, on which date Cassini was de-orbited to burn up in Saturn's upper atmosphere. The Huygens module traveled with Cassini until its separation from the probe on December 25, 2004; Huygens landed by parachute on Titan on January 14, 2005. The separation was facilitated by the SED (Spin/Eject device), which provided a relative separation speed of and a spin rate of 7.5 rpm. It returned data to Earth for around 90 minutes, using the orbiter as a relay. This was the first landing ever accomplished in the outer Solar System and the first landing on a moon other than Earth's Moon. At the end of its mission, the Cassini spacecraft executed its "Grand Finale": a number of risky passes through the gaps between Saturn and its inner rings. This phase aimed to maximize Cassini scientific outcome before the spacecraft was intentionally destroyed to prevent potential contamination of Saturn's moons if Cassini were to unintentionally crash into them when maneuvering the probe was no longer possible due to power loss or other communication issues at the end of its operational lifespan. The atmospheric entry of Cassini ended the mission, but analysis of the returned data will continue for many years. Overview Scientists and individuals from 27 countries made up the joint team responsible for designing, building, flying and collecting data from the Cassini orbiter and the Huygens probe. NASA's Jet Propulsion Laboratory in the United States, where the orbiter was assembled, managed the mission. The European Space Research and Technology Centre developed Huygens. The centre's prime contractor, Aérospatiale of France (part of Thales Alenia Space from 2005), assembled the probe with equipment and instruments supplied by many European countries (including Huygens batteries and two scientific instruments from the United States). The Italian Space Agency (ASI) provided the Cassini orbiter's high-gain radio antenna, with the incorporation of a low-gain antenna (to ensure telecommunications with the Earth for the entire duration of the mission), a compact and lightweight radar, which also used the high-gain antenna and served as a synthetic-aperture radar, a radar altimeter, a radiometer, the radio science subsystem (RSS), and the visible-channel portion VIMS-V of VIMS spectrometer. NASA provided the VIMS infrared counterpart, as well as the Main Electronic Assembly, which included electronic sub-assemblies provided by CNES of France. On April 16, 2008, NASA announced a two-year extension of the funding for ground operations of this mission, at which point it was renamed the Cassini Equinox Mission. The round of funding was again extended in February 2010 with the Cassini Solstice Mission. Naming The mission consisted of two main elements: the ASI/NASA Cassini orbiter, named for the Italian astronomer Giovanni Domenico Cassini, discoverer of Saturn's ring divisions and four of its satellites; and the ESA-developed Huygens probe, named for the Dutch astronomer, mathematician and physicist Christiaan Huygens, discoverer of Titan. The mission was commonly called Saturn Orbiter Titan Probe (SOTP) during gestation, both as a Mariner Mark II mission and generically. Cassini-Huygens was a Flagship-class mission to the outer planets. The other planetary flagships include Galileo, Voyager, and Viking. Objectives Cassini had several objectives, including: Determining the three-dimensional structure and dynamic behavior of the rings of Saturn. Determining the composition of the satellite surfaces and the geological history of each object. Determining the nature and origin of the dark material on Iapetus's leading hemisphere. Measuring the three-dimensional structure and dynamic behavior of the magnetosphere. Studying the dynamic behavior of Saturn's atmosphere at cloud level. Studying the time variability of Titan's clouds and hazes. Characterizing Titan's surface on a regional scale. Cassini–Huygens was launched on October 15, 1997, from Cape Canaveral Air Force Station's Space Launch Complex 40 using a U.S. Air Force Titan IVB/Centaur rocket. The complete launcher was made up of a two-stage Titan IV booster rocket, two strap-on solid rocket engines, the Centaur upper stage, and a payload enclosure, or fairing. The total cost of this scientific exploration mission was about US$3.26 billion, including $1.4 billion for pre-launch development, $704 million for mission operations, $54 million for tracking and $422 million for the launch vehicle. The United States contributed $2.6 billion (80%), the ESA $500 million (15%), and the ASI $160 million (5%). However, these figures are from the press kit which was prepared in October 2000. They do not include inflation over the course of a very long mission, nor do they include the cost of the extended missions. The primary mission for Cassini was completed on July 30, 2008. The mission was extended to June 2010 (Cassini Equinox Mission). This studied the Saturn system in detail during the planet's equinox, which happened in August 2009. On February 3, 2010, NASA announced another extension for Cassini, lasting 6 years until 2017, ending at the time of summer solstice in Saturn's northern hemisphere (Cassini Solstice Mission). The extension enabled another 155 revolutions around the planet, 54 flybys of Titan and 11 flybys of Enceladus. In 2017, an encounter with Titan changed its orbit in such a way that, at closest approach to Saturn, it was only above the planet's cloudtops, below the inner edge of the D ring. This sequence of "proximal orbits" ended when its final encounter with Titan sent the probe into Saturn's atmosphere to be destroyed. Destinations Selected destinations (ordered largest to smallest but not to scale) History Cassini–Huygenss origins date to 1982, when the European Science Foundation and the American National Academy of Sciences formed a working group to investigate future cooperative missions. Two European scientists suggested a paired Saturn Orbiter and Titan Probe as a possible joint mission. In 1983, NASA's Solar System Exploration Committee recommended the same Orbiter and Probe pair as a core NASA project. NASA and the European Space Agency (ESA) performed a joint study of the potential mission from 1984 to 1985. ESA continued with its own study in 1986, while the American astronaut Sally Ride, in her influential 1987 report NASA Leadership and America's Future in Space, also examined and approved of the Cassini mission. While Ride's report described the Saturn orbiter and probe as a NASA solo mission, in 1988 the Associate Administrator for Space Science and Applications of NASA, Len Fisk, returned to the idea of a joint NASA and ESA mission. He wrote to his counterpart at ESA, Roger Bonnet, strongly suggesting that ESA choose the Cassini mission from the three candidates at hand and promising that NASA would commit to the mission as soon as ESA did. At the time, NASA was becoming more sensitive to the strain that had developed between the American and European space programs as a result of European perceptions that NASA had not treated it like an equal during previous collaborations. NASA officials and advisers involved in promoting and planning Cassini–Huygens attempted to correct this trend by stressing their desire to evenly share any scientific and technology benefits resulting from the mission. In part, this newfound spirit of cooperation with Europe was driven by a sense of competition with the Soviet Union, which had begun to cooperate more closely with Europe as ESA drew further away from NASA. Late in 1988, ESA chose Cassini–Huygens as its next major mission and the following year the program received major funding in the US. The collaboration not only improved relations between the two space programs but also helped Cassini–Huygens survive congressional budget cuts in the United States. Cassini–Huygens came under fire politically in both 1992 and 1994, but NASA successfully persuaded the United States Congress that it would be unwise to halt the project after ESA had already poured funds into development because frustration on broken space exploration promises might spill over into other areas of foreign relations. The project proceeded politically smoothly after 1994, although citizens' groups concerned about the potential environmental impact a launch failure might have (because of its plutonium power source) attempted to derail it through protests and lawsuits until and past its 1997 launch. Spacecraft design The spacecraft was planned to be the second three-axis stabilized, RTG-powered Mariner Mark II, a class of spacecraft developed for missions beyond the orbit of Mars, after the Comet Rendezvous Asteroid Flyby (CRAF) mission, but budget cuts and project rescopings forced NASA to terminate CRAF development to save Cassini. As a result, Cassini became more specialized. The Mariner Mark II series was cancelled. The combined orbiter and probe was at the time the third-largest uncrewed interplanetary spacecraft ever successfully launched, behind the Phobos 1 and 2 Mars probes, as well as being among the most complex; NASA's Europa Clipper became the new third-largest probe upon its launch in 2024. The orbiter had a mass of , the probe including of probe support equipment left on the orbiter. With the launch vehicle adapter and of propellants at launch, the spacecraft had a mass of . The Cassini spacecraft was high and wide. Its bus was a dodecagonal prism atop a conical frustum connecting it to a cylinder containing the propellant tanks, to which the RTGs and Huygens were attached. Spacecraft complexity was increased by its trajectory (flight path) to Saturn, and by the ambitious science at its destination. Cassini had 1,630 interconnected electronic components, 22,000 wire connections, and of cabling. The core control computer CPU was a redundant system using the MIL-STD-1750A instruction set architecture. The main propulsion system consisted of one prime and one backup R-4D bipropellant rocket engine. The thrust of each engine was and the total spacecraft delta-v was . Smaller monopropellant rockets provided attitude control. Cassini was powered by of nuclear fuel, mainly plutonium dioxide (containing of pure plutonium). The heat from the material's radioactive decay was turned into electricity. Huygens was supported by Cassini during cruise, but used chemical batteries when independent. The probe contained a DVD with more than 616,400 signatures from citizens in 81 countries, collected in a public campaign. Until September 2017 the Cassini probe continued orbiting Saturn at a distance of between from the Earth. It took 68 to 84 minutes for radio signals to travel from Earth to the spacecraft, and vice versa. Thus ground controllers could not give "real-time" instructions for daily operations or for unexpected events. Even if response were immediate, more than two hours would have passed between the occurrence of a problem and the reception of the engineers' response by the satellite. Instruments Summary Instruments: Optical Remote Sensing ("Located on the remote sensing pallet") Composite Infrared Spectrometer (CIRS) Imaging Science Subsystem (ISS) Ultraviolet Imaging Spectrograph (UVIS) Visible and Infrared Mapping Spectrometer (VIMS) Fields, Particles and Waves (mostly in situ) Cassini Plasma Spectrometer (CAPS) Cosmic Dust Analyzer (CDA) Ion and Neutral Mass Spectrometer (INMS) Magnetometer (MAG) Magnetospheric Imaging Instrument (MIMI) Radio and Plasma Wave Science (RPWS) Microwave Remote Sensing Radar Radio Science (RSS) Description Cassinis instrumentation consisted of: a synthetic aperture radar mapper, a charge-coupled device imaging system, a visible/infrared mapping spectrometer, a composite infrared spectrometer, a cosmic dust analyzer, a radio and plasma wave experiment, a plasma spectrometer, an ultraviolet imaging spectrograph, a magnetospheric imaging instrument, a magnetometer and an ion/neutral mass spectrometer. Telemetry from the communications antenna and other special transmitters (an S-band transmitter and a dual-frequency Ka-band system) was also used to make observations of the atmospheres of Titan and Saturn and to measure the gravity fields of the planet and its satellites. Plutonium power source Because of Saturn's distance from the Sun, solar arrays were not feasible as power sources for this space probe. To generate enough power, such arrays would have been too large and too heavy. Instead, the Cassini orbiter was powered by three GPHS-RTG radioisotope thermoelectric generators, which use heat from the decay of about of plutonium-238 (in the form of plutonium dioxide) to generate direct current electricity via thermoelectrics. The RTGs on the Cassini mission have the same design as those used on the New Horizons, Galileo, and Ulysses space probes, and they were designed to have very long operational lifetimes. At the end of the nominal 11-year Cassini mission, they were still able to produce 600 to 700 watts of electrical power. (Leftover hardware from the Cassini RTG Program was modified and used to power the New Horizons mission to Pluto and the Kuiper belt, which was designed and launched later.) Power distribution was accomplished by 192 solid-state power switches, which also functioned as circuit breakers in the event of an overload condition. The switches used MOSFETs that featured better efficiency and a longer lifetime as compared to conventional switches, while at the same time eliminating transients. However, these solid-state circuit breakers were prone to erroneous tripping (presumably from cosmic rays), requiring them to reset and causing losses in experimental data. To gain momentum while already in flight, the trajectory of the Cassini mission included several gravitational slingshot maneuvers: two fly-by passes of Venus, one more of the Earth, and then one of the planet Jupiter. The terrestrial flyby was the final instance when the probe posed any conceivable danger to human beings. The maneuver was successful, with Cassini passing by above the Earth on August 18, 1999. Had there been any malfunction causing the probe to collide with the Earth, NASA's complete environmental impact study estimated that, in the worst case (with an acute angle of entry in which Cassini would gradually burn up), a significant fraction of the 33 kg of nuclear fuel inside the RTGs would have been dispersed into the Earth's atmosphere so that up to five billion people (i.e. almost the entire terrestrial population) could have been exposed, causing up to an estimated 5,000 additional cancer deaths over the subsequent decades (0.0005 per cent, i.e. a fraction 0.000005, of a billion cancer deaths expected anyway from other causes; the product is incorrectly calculated elsewhere as 500,000 deaths). However, the chance of this happening were estimated to be less than one in one million, i.e. a chance of one person dying (assuming 5,000 deaths) as less than 1 in 200. NASA's risk analysis to use plutonium was publicly criticized by Michio Kaku on the grounds that casualties, property damage, and lawsuits resulting from a possible accident, as well as the potential use of other energy sources, such as solar and fuel cells, were underestimated. Telemetry The Cassini spacecraft was capable of transmitting in several different telemetry formats. The telemetry subsystem is perhaps the most important subsystem, because without it there could be no data return. The telemetry was developed from the ground up, due to the spacecraft using a more modern set of computers than previous missions. Therefore, Cassini was the first spacecraft to adopt mini-packets to reduce the complexity of the Telemetry Dictionary, and the software development process led to the creation of a Telemetry Manager for the mission. There were around 1088 channels (in 67 mini-packets) assembled in the Cassini Telemetry Dictionary. Out of these 67 lower complexity mini-packets, 6 mini-packets contained the subsystem covariance and Kalman gain elements (161 measurements), not used during normal mission operations. This left 947 measurements in 61 mini-packets. A total of seven telemetry maps corresponding to 7 AACS telemetry modes were constructed. These modes are: (1) Record; (2) Nominal Cruise; (3) Medium Slow Cruise; (4) Slow Cruise; (5) Orbital Ops; (6) Av; (7) ATE (Attitude Estimator) Calibration. These 7 maps cover all spacecraft telemetry modes. Huygens probe The Huygens probe, supplied by the European Space Agency (ESA) and named after the 17th century Dutch astronomer who first discovered Titan, Christiaan Huygens, scrutinized the clouds, atmosphere, and surface of Saturn's moon Titan in its descent on January 15, 2005. It was designed to enter and brake in Titan's atmosphere and parachute a fully instrumented robotic laboratory down to the surface. The probe system consisted of the probe itself which descended to Titan, and the probe support equipment (PSE) which remained attached to the orbiting spacecraft. The PSE includes electronics that track the probe, recover the data gathered during its descent, and process and deliver the data to the orbiter that transmits it to Earth. The core control computer CPU was a redundant MIL-STD-1750A control system. The data were transmitted by a radio link between Huygens and Cassini provided by Probe Data Relay Subsystem (PDRS). As the probe's mission could not be telecommanded from Earth because of the great distance, it was automatically managed by the Command Data Management Subsystem (CDMS). The PDRS and CDMS were provided by the Italian Space Agency (ASI). After Cassini launch, it was discovered that data sent from the Huygens probe to Cassini orbiter (and then re-transmitted to Earth) would be largely unreadable. The cause was that the bandwidth of signal processing electronics was too narrow and the anticipated Doppler shift between the lander and the mother craft would put the signals out of the system's range. Thus, Cassini receiver would be unable to receive the data from Huygens during its descent to Titan. A work-around was found to recover the mission. The trajectory of Cassini was altered to reduce the line of sight velocity and therefore the doppler shift. Cassinis subsequent trajectory was identical to the previously planned one, although the change replaced two orbits prior to the Huygens mission with three, shorter orbits. Selected events and discoveries Venus and Earth fly-bys and the cruise to Jupiter The Cassini space probe performed two gravitational-assist flybys of Venus on April 26, 1998, and June 24, 1999. These flybys provided the space probe with enough momentum to travel all the way out to the asteroid belt, while the Sun's gravity pulled the space probe back into the inner Solar System. On August 18, 1999, at 03:28 UTC, the craft made a gravitational-assist flyby of the Earth. One hour and 20 minutes before closest approach, Cassini made its closest approach to the Earth's Moon at 377,000 kilometers, and it took a series of calibration photos. On January 23, 2000, Cassini performed a flyby of the asteroid 2685 Masursky at around 10:00 UTC. It took photos in the period five to seven hours before the flyby at a distance of and a diameter of was estimated for the asteroid. Jupiter flyby Cassini made its closest approach to Jupiter on December 30, 2000, at 9.7 million kilometers, and made many scientific measurements. About 26,000 images of Jupiter, its faint rings, and its moons were taken during the six-month flyby. It produced the most detailed global color portrait of the planet yet (see image at right), in which the smallest visible features are approximately across. A major finding of the flyby, announced on March 6, 2003, was of Jupiter's atmospheric circulation. Dark "belts" alternate with light "zones" in the atmosphere, and scientists had long considered the zones, with their pale clouds, to be areas of upwelling air, partly because many clouds on Earth form where air is rising. But analysis of Cassini imagery showed that individual storm cells of upwelling bright-white clouds, too small to see from Earth, pop up almost without exception in the dark belts. According to Anthony Del Genio of NASA's Goddard Institute for Space Studies, "the belts must be the areas of net-rising atmospheric motion on Jupiter, [so] the net motion in the zones has to be sinking". Other atmospheric observations included a swirling dark oval of high atmospheric haze, about the size of the Great Red Spot, near Jupiter's north pole. Infrared imagery revealed aspects of circulation near the poles, with bands of globe-encircling winds, with adjacent bands moving in opposite directions. The same announcement also discussed the nature of Jupiter's rings. Light scattering by particles in the rings showed the particles were irregularly shaped (rather than spherical) and likely originate as ejecta from micrometeorite impacts on Jupiter's moons, probably Metis and Adrastea. Tests of general relativity On October 10, 2003, the mission's science team announced the results of tests of Albert Einstein's general theory of relativity, performed by using radio waves transmitted from the Cassini space probe. The radio scientists measured a frequency shift in the radio waves to and from the spacecraft, as they passed close to the Sun. According to the general theory of relativity, a massive object like the Sun causes space-time to curve, causing a beam of radiowaves travelling out of its gravitational well to decrease in frequency and radiowaves travelling into the gravitational well to increase in frequency, referred to as gravitational redshift / blueshift. Although some measurable deviations from the values calculated using the general theory of relativity are predicted by some unusual cosmological models, no such deviations were found by this experiment. Previous tests using radiowaves transmitted by the Viking and Voyager space probes were in agreement with the calculated values from general relativity to within an accuracy of one part in one thousand. The more refined measurements from the Cassini space probe experiment improved this accuracy to about one part in 51,000. The data firmly support Einstein's general theory of relativity. New moons of Saturn In total, the Cassini mission discovered seven new moons orbiting Saturn. Using images taken by Cassini, researchers discovered Methone, Pallene and Polydeuces in 2004, although later analysis revealed that Voyager 2 had photographed Pallene in its 1981 flyby of the ringed planet. On May 1, 2005, a new moon was discovered by Cassini in the Keeler gap. It was given the designation S/2005 S 1 before being named Daphnis. A fifth new moon was discovered by Cassini on May 30, 2007, and was provisionally labeled S/2007 S 4. It is now known as Anthe. A press release on February 3, 2009, showed a sixth new moon found by Cassini. The moon is approximately in diameter within the G-ring of the ring system of Saturn, and is now named Aegaeon (formerly S/2008 S 1). A press release on November 2, 2009, mentions the seventh new moon found by Cassini on July 26, 2009. It is presently labeled S/2009 S 1 and is approximately in diameter in the B-ring system. On April 14, 2014, NASA scientists reported the possible beginning of a new moon in Saturn's A Ring. Phoebe flyby On June 11, 2004, Cassini flew by the moon Phoebe. This was the first opportunity for close-up studies of this moon (Voyager 2 performed a distant flyby in 1981 but returned no detailed images). It also was Cassini's only possible flyby for Phoebe due to the mechanics of the available orbits around Saturn. The first close-up images were received on June 12, 2004, and mission scientists immediately realized that the surface of Phoebe looks different from asteroids visited by spacecraft. Parts of the heavily cratered surface look very bright in those pictures, and it is currently believed that a large amount of water ice exists under its immediate surface. Saturn rotation In an announcement on June 28, 2004, Cassini program scientists described the measurement of the rotational period of Saturn. Because there are no fixed features on the surface that can be used to obtain this period, the repetition of radio emissions was used. This new data agreed with the latest values measured from Earth, and constituted a puzzle to the scientists. It turns out that the radio rotational period had changed since it was first measured in 1980 by Voyager 1, and it was now 6 minutes longer. This, however, does not indicate a change in the overall spin of the planet. It is thought to be due to variations in the upper atmosphere and ionosphere at the latitudes which are magnetically connected to the radio source region. In 2019 NASA announced Saturn's rotational period as 10 hours, 33 minutes, 38 seconds, calculated using Saturnian ring seismology. Vibrations from Saturn's interior cause oscillations in its gravitational field. This energy is absorbed by ring particles in specific locations, where it accumulates until it is released in a wave. Scientists used data from more than 20 of these waves to construct a family of models of Saturn's interior, providing basis for calculating its rotational period. Orbiting Saturn On July 1, 2004, the spacecraft flew through the gap between the F and G rings and achieved orbit, after a seven-year voyage. It was the first spacecraft to orbit Saturn. The Saturn Orbital Insertion (SOI) maneuver performed by Cassini was complex, requiring the craft to orient its High-Gain Antenna away from Earth and along its flight path, to shield its instruments from particles in Saturn's rings. Once the craft crossed the ring plane, it had to rotate again to point its engine along its flight path, and then the engine fired to decelerate the craft by 622 m/s to allow Saturn to capture it. Cassini was captured by Saturn's gravity at around 8:54 pm Pacific Daylight Time on June 30, 2004. During the maneuver Cassini passed within of Saturn's cloud tops. When Cassini was in Saturnian orbit, departure from the Saturn system was evaluated in 2008 during end of mission planning. Titan flybys Cassini had its first flyby of Saturn's largest moon, Titan, on July 2, 2004, a day after orbit insertion, when it approached to within of Titan. Images taken through special filters (able to see through the moon's global haze) showed south polar clouds thought to be composed of methane and surface features with widely differing brightness. On October 27, 2004, the spacecraft executed the first of the 45 planned close flybys of Titan when it passed a mere above the moon. Almost four gigabits of data were collected and transmitted to Earth, including the first radar images of the moon's haze-enshrouded surface. It revealed the surface of Titan (at least the area covered by radar) to be relatively level, with topography reaching no more than about in altitude. The flyby provided a remarkable increase in imaging resolution over previous coverage. Images with up to 100 times better resolution were taken and are typical of resolutions planned for subsequent Titan flybys. Cassini collected pictures of Titan and the lakes of methane were similar to the lakes of water on Earth. Huygens lands on Titan Cassini released the Huygens probe on December 25, 2004, by means of a spring and spiral rails intended to rotate the probe for greater stability. It entered the atmosphere of Titan on January 14, 2005, and after a two-and-a-half-hour descent landed on solid ground. Although Cassini successfully relayed 350 of the pictures that it received from Huygens of its descent and landing site, a malfunction in one of the communications channels resulted in the loss of a further 350 pictures. Enceladus flybys During the first two close flybys of the moon Enceladus in 2005, Cassini discovered a deflection in the local magnetic field that is characteristic for the existence of a thin but significant atmosphere. Other measurements obtained at that time point to ionized water vapor as its main constituent. Cassini also observed water ice geysers erupting from the south pole of Enceladus, which gives more credibility to the idea that Enceladus is supplying the particles of Saturn's E ring. Mission scientists began to suspect that there may be pockets of liquid water near the surface of the moon that fuel the eruptions. On March 12, 2008, Cassini made a close fly-by of Enceladus, passing within 50 km of the moon's surface. The spacecraft passed through the plumes extending from its southern geysers, detecting water, carbon dioxide and various hydrocarbons with its mass spectrometer, while also mapping surface features that are at much higher temperature than their surroundings with the infrared spectrometer. Cassini was unable to collect data with its cosmic dust analyzer due to an unknown software malfunction. On November 21, 2009, Cassini made its eighth flyby of Enceladus, this time with a different geometry, approaching within of the surface. The Composite Infrared Spectrograph (CIRS) instrument produced a map of thermal emissions from the Baghdad Sulcus 'tiger stripe'. The data returned helped create a detailed and high resolution mosaic image of the southern part of the moon's Saturn-facing hemisphere. On April 3, 2014, nearly ten years after Cassini entered Saturn's orbit, NASA reported evidence of a large salty internal ocean of liquid water in Enceladus. The presence of an internal salty ocean in contact with the moon's rocky core, places Enceladus "among the most likely places in the Solar System to host alien microbial life". On June 30, 2014, NASA celebrated ten years of Cassini exploring Saturn and its moons, highlighting the discovery of water activity on Enceladus among other findings. In September 2015, NASA announced that gravitational and imaging data from Cassini were used to analyze the librations of Enceladus' orbit and determined that the moon's surface is not rigidly joined to its core, concluding that the underground ocean must therefore be global in extent. On October 28, 2015, Cassini performed a close flyby of Enceladus, coming within of the surface, and passing through the icy plume above the south pole. On December 14, 2023, astronomers reported the first time discovery, in the plumes of Enceladus, of hydrogen cyanide, a possible chemical essential for life as we know it, as well as other organic molecules, some of which are yet to be better identified and understood. According to the researchers, "these [newly discovered] compounds could potentially support extant microbial communities or drive complex organic synthesis leading to the origin of life". Radio occultations of Saturn's rings In May 2005, Cassini began a series of radio occultation experiments, to measure the size-distribution of particles in Saturn's rings, and measure the atmosphere of Saturn itself. For over four months, the craft completed orbits designed for this purpose. During these experiments, it flew behind the ring plane of Saturn, as seen from Earth, and transmitted radio waves through the particles. The radio signals received on Earth were analyzed, for frequency, phase, and power shift of the signal to determine the structure of the rings. Spokes in rings verified In images captured September 5, 2005, Cassini detected spokes in Saturn's rings, previously seen only by the visual observer Stephen James O'Meara in 1977 and then confirmed by the Voyager space probes in the early 1980s. Lakes of Titan Radar images obtained on July 21, 2006, appear to show lakes of liquid hydrocarbon (such as methane and ethane) in Titan's northern latitudes. This is the first discovery of currently existing lakes anywhere besides on Earth. The lakes range in size from one to one-hundred kilometers across. On March 13, 2007, the Jet Propulsion Laboratory announced that it had found strong evidence of seas of methane and ethane in the northern hemisphere of Titan. At least one of these is larger than any of the Great Lakes in North America. Saturn hurricane In November 2006, scientists discovered a storm at the south pole of Saturn with a distinct eyewall. This is characteristic of a hurricane on Earth and had never been seen on another planet before. Unlike a terrestrial hurricane, the storm appears to be stationary at the pole. The storm is across, and high, with winds blowing at . Iapetus flyby On September 10, 2007, Cassini completed its flyby of the strange, two-toned, walnut-shaped moon, Iapetus. Images were taken from above the surface. As it was sending the images back to Earth, it was hit by a cosmic ray that forced it to temporarily enter safe mode. All of the data from the flyby were recovered. Mission extension On April 15, 2008, Cassini received funding for a 27-month extended mission. It consisted of 60 more orbits of Saturn, with 21 more close Titan flybys, seven of Enceladus, six of Mimas, eight of Tethys, and one targeted flyby each of Dione, Rhea, and Helene. The extended mission began on July 1, 2008, and was renamed the Cassini Equinox Mission as the mission coincided with Saturn's equinox. Second mission extension A proposal was submitted to NASA for a second mission extension (September 2010 – May 2017), provisionally named the extended-extended mission or XXM. This ($60M pa) was approved in February 2010 and renamed the Cassini Solstice Mission. It included Cassini orbiting Saturn 155 more times, conducting 54 additional flybys of Titan and 11 more of Enceladus. Great Storm of 2010 and aftermath On October 25, 2012, Cassini witnessed the aftermath of the massive Great White Spot storm that recurs roughly every 30 years on Saturn. Data from the composite infrared spectrometer (CIRS) instrument indicated a powerful discharge from the storm that caused a temperature spike in the stratosphere of Saturn above normal. Simultaneously, a huge increase in ethylene gas was detected by NASA researchers at Goddard Research Center in Greenbelt, Maryland. Ethylene is a colorless gas that is highly uncommon on Saturn and is produced both naturally and through man-made sources on Earth. The storm that produced this discharge was first observed by the spacecraft on December 5, 2010, in Saturn's northern hemisphere. The storm is the first of its kind to be observed by a spacecraft in orbit around Saturn as well as the first to be observed at thermal infrared wavelengths, allowing scientists to observe the temperature of Saturn's atmosphere and track phenomena that are invisible to the naked eye. The spike of ethylene gas that was produced by the storm reached levels that were 100 times more than those thought possible for Saturn. Scientists have also determined that the storm witnessed was the largest, hottest stratospheric vortex ever detected in the Solar System, initially being larger than Jupiter's Great Red Spot. Venus transit On December 21, 2012, Cassini observed a transit of Venus across the Sun. The VIMS instrument analyzed sunlight passing through the Venusian atmosphere. VIMS previously observed the transit of exoplanet HD 189733 b. The Day the Earth Smiled On July 19, 2013, the probe was pointed towards Earth to capture an image of the Earth and the Moon, as part of a natural light, multi-image portrait of the entire Saturn system. The event was unique as it was the first time NASA informed the public that a long-distance photo was being taken in advance. The imaging team said they wanted people to smile and wave to the skies, with Cassini scientist Carolyn Porco describing the moment as a chance to "celebrate life on the Pale Blue Dot". Rhea flyby On February 10, 2015, the Cassini spacecraft visited Rhea more closely, coming within . The spacecraft observed the moon with its cameras producing some of the highest resolution color images yet of Rhea. Hyperion flyby Cassini performed its latest flyby of Saturn's moon Hyperion on May 31, 2015, at a distance of about . Dione flyby Cassini performed its last flyby of Saturn's moon Dione on August 17, 2015, at a distance of about . A previous flyby was performed on June 16. Hexagon changes color Between 2012 and 2016, the persistent hexagonal cloud pattern at Saturn's north pole changed from a mostly blue color to more of a golden color. One theory for this is a seasonal change: extended exposure to sunlight may be creating haze as the pole swivels toward the Sun. It was previously noted that there was less blue color overall on Saturn between 2004 and 2008. Grand Finale and destruction Cassini end involved a series of close Saturn passes, approaching within the rings, then an entry into Saturn's atmosphere on September 15, 2017, to destroy the spacecraft. This method was chosen to ensure protection and prevent biological contamination to any of the moons of Saturn thought to offer potential habitability. In 2008 a number of options were evaluated to achieve this goal, each with varying funding, scientific, and technical challenges. A short period Saturn impact for an end of mission was rated "excellent" with the reasons "D-ring option satisfies unachieved AO goals; cheap and easily achievable" while collision with an icy moon was rated "good" for being "cheap and achievable anywhere/time". There were problems in 2013–14 about NASA receiving U.S. government funding for the Grand Finale. The two phases of the Grand Finale ended up being the equivalent of having two separate Discovery-class missions in that the Grand Finale was completely different from the main Cassini regular mission. The U.S. government in late 2014 approved the Grand Finale at the cost of $200 million. This was far cheaper than building two new probes in separate Discovery-class missions. On November 29, 2016, the spacecraft performed a Titan flyby that took it to the gateway of F-ring orbits: This was the start of the Grand Finale phase culminating in its impact with the planet. A final Titan flyby on April 22, 2017, changed the orbit again to fly through the gap between Saturn and its inner ring days later on April 26. Cassini passed about above Saturn's cloud layer and from the visible edge of the inner ring; it successfully took images of Saturn's atmosphere and began returning data the next day. After a further 22 orbits through the gap, the mission was ended with a dive into Saturn's atmosphere on September 15; signal was lost at 11:55:46 UTC on September 15, 2017, just 30 seconds later than predicted. It is estimated that the spacecraft burned up about 45 seconds after the last transmission. In September 2018, NASA won an Emmy Award for Outstanding Original Interactive Program for its presentation of the Cassini mission's Grand Finale at Saturn. In December 2018, Netflix aired "NASA's Cassini Mission" on their series 7 Days Out documenting the final days of work on the Cassini mission before the spacecraft crashed into Saturn to complete its Grand Finale. In January 2019, new research using data collected during Cassini Grand Finale phase was published: The final close passes by the rings and planet enabled scientists to measure the length of a day on Saturn: 10 hours, 33 minutes and 38 seconds. Saturn's rings are relatively new, 10 to 100 million years old. Missions The spacecraft operation was organized around a series of missions. Each is structured according to a certain amount of funding, goals, etc. At least 260 scientists from 17 countries have worked on the Cassini–Huygens mission; in addition thousands of people overall worked to design, manufacture, and launch the mission. Prime Mission, July 2004 through June 2008. Cassini Equinox Mission was a two-year mission extension which ran from July 2008 through September 2010. Cassini Solstice Mission ran from October 2010 through April 2017. (Also known as the XXM mission.) Grand Finale (spacecraft directed into Saturn), April 2017 to September 15, 2017. Glossary AACS: Attitude and Articulation Control Subsystem ACS: Attitude Control Subsystem AFC: AACS Flight Computer ARWM: Articulated Reaction Wheel Mechanism ASI: Agenzia Spaziale Italiana, the Italian space agency BIU: Bus Interface Unit BOL: Beginning of Life CAM: Command Approval Meeting CDS: Command and Data Subsystem—Cassini computer that commands and collects data from the instruments CICLOPS: Cassini Imaging Central Laboratory for Operations CIMS: Cassini Information Management System CIRS: Composite Infrared Spectrometer DCSS: Descent Control Subsystem DSCC: Deep Space Communications Center DSN: Deep Space Network (large antennas around the Earth) DTSTART: Dead Time Start ELS: Electron Spectrometer (part of CAPS instrument) EOM: End of Mission ERT: Earth-received time, UTC of an event ESA: European Space Agency ESOC: European Space Operations Centre FSW: flight software HGA: High Gain Antenna HMCS: Huygens Monitoring and Control System HPOC: Huygens Probe Operations Center IBS: Ion Beam Spectrometer (part of CAPS instrument) IEB: Instrument Expanded Blocks (instrument command sequences) IMS: Ion Mass Spectrometer (part of CAPS instrument) ITL: Integrated Test Laboratory—spacecraft simulator IVP: Inertial Vector Propagator LGA: Low Gain Antenna NAC: Narrow Angle Camera NASA: National Aeronautics and Space Administration, the United States space agency OTM: Orbit Trim Maneuver PDRS: Probe Data Relay Subsystem PHSS: Probe Harness SubSystem POSW: Probe On-Board Software PPS: Power and Pyrotechnic Subsystem PRA: Probe Relay Antenna PSA: Probe Support Avionics PSIV: Preliminary Sequence Integration and Validation PSE: probe support equipment RCS: Reaction Control System RFS: Radio Frequency Subsystem RPX: ring plane crossing RWA: Reaction Wheel Assembly SCET: Spacecraft Event Time SCR: sequence change requests SKR: Saturn Kilometric Radiation SOI: Saturn Orbit Insertion (July 1, 2004) SOP: Science Operations Plan SSPS: Solid State Power Switch SSR: Solid State Recorder SSUP: Science and Sequence Update Process TLA: Thermal Louver Assemblies USO: UltraStable Oscillator VRHU: Variable Radioisotope Heater Units WAC: Wide Angle Camera XXM: Extended-Extended Mission	Technology	Unmanned spacecraft	null
67958	https://en.wikipedia.org/wiki/Copernicium	Copernicium	Copernicium is a synthetic chemical element; it has symbol Cn and atomic number 112. Its known isotopes are extremely radioactive, and have only been created in a laboratory. The most stable known isotope, copernicium-285, has a half-life of approximately 30 seconds. Copernicium was first created in February 1996 by the GSI Helmholtz Centre for Heavy Ion Research near Darmstadt, Germany. It was named after the astronomer Nicolaus Copernicus on his 537th anniversary. In the periodic table of the elements, copernicium is a d-block transactinide element and a group 12 element. During reactions with gold, it has been shown to be an extremely volatile element, so much so that it is possibly a gas or a volatile liquid at standard temperature and pressure. Copernicium is calculated to have several properties that differ from its lighter homologues in group 12, zinc, cadmium and mercury; due to relativistic effects, it may give up its 6d electrons instead of its 7s ones, and it may have more similarities to the noble gases such as radon rather than its group 12 homologues. Calculations indicate that copernicium may show the oxidation state +4, while mercury shows it in only one compound of disputed existence and zinc and cadmium do not show it at all. It has also been predicted to be more difficult to oxidize copernicium from its neutral state than the other group 12 elements. Predictions vary on whether solid copernicium would be a metal, semiconductor, or insulator. Copernicium is one of the heaviest elements whose chemical properties have been experimentally investigated. Introduction History Discovery Copernicium was first created on February 9, 1996, at the Gesellschaft für Schwerionenforschung (GSI) in Darmstadt, Germany, by Sigurd Hofmann, Victor Ninov et al. This element was created by firing accelerated zinc-70 nuclei at a target made of lead-208 nuclei in a heavy ion accelerator. A single atom of copernicium was produced with a mass number of 277. (A second was originally reported, but was found to have been based on data fabricated by Ninov, and was thus retracted.) Pb + Zn → Cn* → Cn + n In May 2000, the GSI successfully repeated the experiment to synthesize a further atom of copernicium-277. This reaction was repeated at RIKEN using the Search for a Super-Heavy Element Using a Gas-Filled Recoil Separator set-up in 2004 and 2013 to synthesize three further atoms and confirm the decay data reported by the GSI team. This reaction had also previously been tried in 1971 at the Joint Institute for Nuclear Research in Dubna, Russia to aim for 276Cn (produced in the 2n channel), but without success. The IUPAC/IUPAP Joint Working Party (JWP) assessed the claim of copernicium's discovery by the GSI team in 2001 and 2003. In both cases, they found that there was insufficient evidence to support their claim. This was primarily related to the contradicting decay data for the known nuclide rutherfordium-261. However, between 2001 and 2005, the GSI team studied the reaction 248Cm(26Mg,5n)269Hs, and were able to confirm the decay data for hassium-269 and rutherfordium-261. It was found that the existing data on rutherfordium-261 was for an isomer, now designated rutherfordium-261m. In May 2009, the JWP reported on the claims of discovery of element 112 again and officially recognized the GSI team as the discoverers of element 112. This decision was based on the confirmation of the decay properties of daughter nuclei as well as the confirmatory experiments at RIKEN. Work had also been done at the Joint Institute for Nuclear Research in Dubna, Russia from 1998 to synthesise the heavier isotope 283Cn in the hot fusion reaction 238U(48Ca,3n)283Cn; most observed atoms of 283Cn decayed by spontaneous fission, although an alpha decay branch to 279Ds was detected. While initial experiments aimed to assign the produced nuclide with its observed long half-life of 3 minutes based on its chemical behaviour, this was found to be not mercury-like as would have been expected (copernicium being under mercury in the periodic table), and indeed now it appears that the long-lived activity might not have been from 283Cn at all, but its electron capture daughter 283Rg instead, with a shorter 4-second half-life associated with 283Cn. (Another possibility is assignment to a metastable isomeric state, 283mCn.) While later cross-bombardments in the 242Pu+48Ca and 245Cm+48Ca reactions succeeded in confirming the properties of 283Cn and its parents 287Fl and 291Lv, and played a major role in the acceptance of the discoveries of flerovium and livermorium (elements 114 and 116) by the JWP in 2011, this work originated subsequent to the GSI's work on 277Cn and priority was assigned to the GSI. Naming Using Mendeleev's nomenclature for unnamed and undiscovered elements, copernicium should be known as eka-mercury. In 1979, IUPAC published recommendations according to which the element was to be called ununbium (with the corresponding symbol of Uub), a systematic element name as a placeholder, until the element was discovered (and the discovery then confirmed) and a permanent name was decided on. Although widely used in the chemical community on all levels, from chemistry classrooms to advanced textbooks, the recommendations were mostly ignored among scientists in the field, who either called it "element 112", with the symbol of E112, (112), or even simply 112. After acknowledging the GSI team's discovery, the IUPAC asked them to suggest a permanent name for element 112. On 14 July 2009, they proposed copernicium with the element symbol Cp, after Nicolaus Copernicus "to honor an outstanding scientist, who changed our view of the world". During the standard six-month discussion period among the scientific community about the naming, it was pointed out that the symbol Cp was previously associated with the name cassiopeium (cassiopium), now known as lutetium (Lu). Moreover, Cp is frequently used today to mean the cyclopentadienyl ligand (C5H5). Primarily because cassiopeium (Cp) was (until 1949) accepted by IUPAC as an alternative allowed name for lutetium, the IUPAC disallowed the use of Cp as a future symbol, prompting the GSI team to put forward the symbol Cn as an alternative. On 19 February 2010, the 537th anniversary of Copernicus' birth, IUPAC officially accepted the proposed name and symbol. Isotopes Copernicium has no stable or naturally occurring isotopes. Several radioactive isotopes have been synthesized in the laboratory, either by fusing two atoms or by observing the decay of heavier elements. Eight different isotopes have been reported with mass numbers 277 and 280–286, and one unconfirmed metastable isomer in 285Cn has been reported. Most of these decay predominantly through alpha decay, but some undergo spontaneous fission, and copernicium-283 may have an electron capture branch. The isotope copernicium-283 was instrumental in the confirmation of the discoveries of the elements flerovium and livermorium. Half-lives All confirmed copernicium isotopes are extremely unstable and radioactive; in general, heavier isotopes are more stable than the lighter, and isotopes with an odd neutron number have relatively longer half-lives due to additional hindrance against spontaneous fission. The most stable known isotope, 285Cn, has a half-life of 30 seconds; 283Cn has a half-life of 4 seconds, and the unconfirmed 285mCn and 286Cn have half-lives of about 15 and 8.45 seconds respectively. Other isotopes have half-lives shorter than one second. 281Cn and 284Cn both have half-lives on the order of 0.1 seconds, and the remaining isotopes have half-lives shorter than one millisecond. It is predicted that the heavy isotopes 291Cn and 293Cn may have half-lives longer than a few decades, for they are predicted to lie near the center of the theoretical island of stability, and may have been produced in the r-process and be detectable in cosmic rays, though they would be about 10−12 times as abundant as lead. The lightest isotopes of copernicium have been synthesized by direct fusion between two lighter nuclei and as decay products (except for 277Cn, which is not known to be a decay product), while the heavier isotopes are only known to be produced by decay of heavier nuclei. The heaviest isotope produced by direct fusion is 283Cn; the three heavier isotopes, 284Cn, 285Cn, and 286Cn, have only been observed as decay products of elements with larger atomic numbers. In 1999, American scientists at the University of California, Berkeley, announced that they had succeeded in synthesizing three atoms of 293Og. These parent nuclei were reported to have successively emitted three alpha particles to form copernicium-281 nuclei, which were claimed to have undergone alpha decay, emitting alpha particles with decay energy 10.68 MeV and half-life 0.90 ms, but their claim was retracted in 2001 as it had been based on data fabricated by Ninov. This isotope was truly produced in 2010 by the same team; the new data contradicted the previous fabricated data. The missing isotopes 278Cn and 279Cn are too heavy to be produced by cold fusion and too light to be produced by hot fusion. They might be filled from above by decay of heavier elements produced by hot fusion, and indeed 280Cn and 281Cn were produced this way. The isotopes 286Cn and 287Cn could be produced by charged-particle evaporation, in the reaction 244Pu(48Ca,αxn) with x equalling 1 or 2. Predicted properties Very few properties of copernicium or its compounds have been measured; this is due to its extremely limited and expensive production and the fact that copernicium (and its parents) decays very quickly. A few singular chemical properties have been measured, as well as the boiling point, but properties of the copernicium metal remain generally unknown and for the most part, only predictions are available. Chemical Copernicium is the tenth and last member of the 6d series and is the heaviest group 12 element in the periodic table, below zinc, cadmium and mercury. It is predicted to differ significantly from the lighter group 12 elements. The valence s-subshells of the group 12 elements and period 7 elements are expected to be relativistically contracted most strongly at copernicium. This and the closed-shell configuration of copernicium result in it probably being a very noble metal. A standard reduction potential of +2.1 V is predicted for the Cn2+/Cn couple. Copernicium's predicted first ionization energy of 1155 kJ/mol almost matches that of the noble gas xenon at 1170.4 kJ/mol. Copernicium's metallic bonds should also be very weak, possibly making it extremely volatile like the noble gases, and potentially making it gaseous at room temperature. However, it should be able to form metal–metal bonds with copper, palladium, platinum, silver, and gold; these bonds are predicted to be only about 15–20 kJ/mol weaker than the analogous bonds with mercury. In opposition to the earlier suggestion, ab initio calculations at the high level of accuracy predicted that the chemistry of singly-valent copernicium resembles that of mercury rather than that of the noble gases. The latter result can be explained by the huge spin–orbit interaction which significantly lowers the energy of the vacant 7p1/2 state of copernicium. Once copernicium is ionized, its chemistry may present several differences from those of zinc, cadmium, and mercury. Due to the stabilization of 7s electronic orbitals and destabilization of 6d ones caused by relativistic effects, Cn2+ is likely to have a [Rn]5f146d87s2 electronic configuration, using the 6d orbitals before the 7s one, unlike its homologues. The fact that the 6d electrons participate more readily in chemical bonding means that once copernicium is ionized, it may behave more like a transition metal than its lighter homologues, especially in the possible +4 oxidation state. In aqueous solutions, copernicium may form the +2 and perhaps +4 oxidation states. The diatomic ion , featuring mercury in the +1 oxidation state, is well-known, but the ion is predicted to be unstable or even non-existent. Copernicium(II) fluoride, CnF2, should be more unstable than the analogous mercury compound, mercury(II) fluoride (HgF2), and may even decompose spontaneously into its constituent elements. As the most electronegative reactive element, fluorine may be the only element able to oxidise copernicium even further to the +4 and even +6 oxidation states in CnF4 and CnF6; the latter may require matrix-isolation conditions to be detected, as in the disputed detection of HgF4. CnF4 should be more stable than CnF2. In polar solvents, copernicium is predicted to preferentially form the and anions rather than the analogous neutral fluorides (CnF4 and CnF2, respectively), although the analogous bromide or iodide ions may be more stable towards hydrolysis in aqueous solution. The anions and should also be able to exist in aqueous solution. The formation of thermodynamically stable copernicium(II) and (IV) fluorides would be analogous to the chemistry of xenon. Analogous to mercury(II) cyanide (Hg(CN)2), copernicium is expected to form a stable cyanide, Cn(CN)2. Physical and atomic Copernicium should be a dense metal, with a density of 14.0 g/cm3 in the liquid state at 300 K; this is similar to the known density of mercury, which is 13.534 g/cm3. (Solid copernicium at the same temperature should have a higher density of 14.7 g/cm3.) This results from the effects of copernicium's higher atomic weight being cancelled out by its larger interatomic distances compared to mercury. Some calculations predicted copernicium to be a gas at room temperature due to its closed-shell electron configuration, which would make it the first gaseous metal in the periodic table. A 2019 calculation agrees with these predictions on the role of relativistic effects, suggesting that copernicium will be a volatile liquid bound by dispersion forces under standard conditions. Its melting point is estimated at and its boiling point at , the latter in agreement with the experimentally estimated value of . The atomic radius of copernicium is expected to be around 147 pm. Due to the relativistic stabilization of the 7s orbital and destabilization of the 6d orbital, the Cn+ and Cn2+ ions are predicted to give up 6d electrons instead of 7s electrons, which is the opposite of the behavior of its lighter homologues. In addition to the relativistic contraction and binding of the 7s subshell, the 6d5/2 orbital is expected to be destabilized due to spin–orbit coupling, making it behave similarly to the 7s orbital in terms of size, shape, and energy. Predictions of the expected band structure of copernicium are varied. Calculations in 2007 expected that copernicium may be a semiconductor with a band gap of around 0.2 eV, crystallizing in the hexagonal close-packed crystal structure. However, calculations in 2017 and 2018 suggested that copernicium should be a noble metal at standard conditions with a body-centered cubic crystal structure: it should hence have no band gap, like mercury, although the density of states at the Fermi level is expected to be lower for copernicium than for mercury. 2019 calculations then suggested that in fact copernicium has a large band gap of 6.4 ± 0.2 eV, which should be similar to that of the noble gas radon (predicted as 7.1 eV) and would make it an insulator; bulk copernicium is predicted by these calculations to be bound mostly by dispersion forces, like the noble gases. Like mercury, radon, and flerovium, but not oganesson (eka-radon), copernicium is calculated to have no electron affinity. Experimental atomic gas phase chemistry Interest in copernicium's chemistry was sparked by predictions that it would have the largest relativistic effects in the whole of period 7 and group 12, and indeed among all 118 known elements. Copernicium is expected to have the ground state electron configuration [Rn] 5f14 6d10 7s2 and thus should belong to group 12 of the periodic table, according to the Aufbau principle. As such, it should behave as the heavier homologue of mercury and form strong binary compounds with noble metals like gold. Experiments probing the reactivity of copernicium have focused on the adsorption of atoms of element 112 onto a gold surface held at varying temperatures, in order to calculate an adsorption enthalpy. Owing to relativistic stabilization of the 7s electrons, copernicium shows radon-like properties. Experiments were performed with the simultaneous formation of mercury and radon radioisotopes, allowing a comparison of adsorption characteristics. The first chemical experiments on copernicium were conducted using the 238U(48Ca,3n)283Cn reaction. Detection was by spontaneous fission of the claimed parent isotope with half-life of 5 minutes. Analysis of the data indicated that copernicium was more volatile than mercury and had noble gas properties. However, the confusion regarding the synthesis of copernicium-283 has cast some doubt on these experimental results. Given this uncertainty, between April–May 2006 at the JINR, a FLNR–PSI team conducted experiments probing the synthesis of this isotope as a daughter in the nuclear reaction 242Pu(48Ca,3n)287Fl. (The 242Pu + 48Ca fusion reaction has a slightly larger cross-section than the 238U + 48Ca reaction, so that the best way to produce copernicium for chemical experimentation is as an overshoot product as the daughter of flerovium.) In this experiment, two atoms of copernicium-283 were unambiguously identified and the adsorption properties were interpreted to show that copernicium is a more volatile homologue of mercury, due to formation of a weak metal-metal bond with gold. This agrees with general indications from some relativistic calculations that copernicium is "more or less" homologous to mercury. However, it was pointed out in 2019 that this result may simply be due to strong dispersion interactions. In April 2007, this experiment was repeated and a further three atoms of copernicium-283 were positively identified. The adsorption property was confirmed and indicated that copernicium has adsorption properties in agreement with being the heaviest member of group 12. These experiments also allowed the first experimental estimation of copernicium's boiling point: 84 °C, so that it may be a gas at standard conditions. Because the lighter group 12 elements often occur as chalcogenide ores, experiments were conducted in 2015 to deposit copernicium atoms on a selenium surface to form copernicium selenide, CnSe. Reaction of copernicium atoms with trigonal selenium to form a selenide was observed, with -ΔHadsCn(t-Se) > 48 kJ/mol, with the kinetic hindrance towards selenide formation being lower for copernicium than for mercury. This was unexpected as the stability of the group 12 selenides tends to decrease down the group from ZnSe to HgSe.	Physical sciences	Group 12	Chemistry
67963	https://en.wikipedia.org/wiki/Ginkgo%20biloba	Ginkgo biloba	Ginkgo biloba, commonly known as ginkgo or gingko ( ), also known as the maidenhair tree, is a species of gymnosperm tree native to East Asia. It is the last living species in the order Ginkgoales, which first appeared over 290 million years ago, and fossils very similar to the living species, belonging to the genus Ginkgo, extend back to the Middle Jurassic epoch approximately 170 million years ago. The tree was cultivated early in human history and remains commonly planted, and is widely regarded as a living fossil. The plant may be toxic or allergenic in certain cases. Leaf extract is commonly used as a dietary supplement, but there is insufficient clinical evidence that it supports human health or is effective against any disease. Description Ginkgos are large trees, normally reaching a height of , with some specimens in China being over . The tree has an angular crown and long, somewhat erratic branches, and is usually deep-rooted and resistant to wind and snow damage. Young trees are often tall and slender, and sparsely branched; the crown becomes broader as the tree ages. A combination of resistance to disease, insect-resistant wood, and the ability to form aerial roots and sprouts makes ginkgos durable, with some specimens claimed to be more than 2,500 years old. Leaves The leaves are unique among seed plants, being fan-shaped with veins radiating out into the leaf blade, sometimes bifurcating (splitting), but never anastomosing to form a network. Two veins enter the leaf blade at the base and fork repeatedly in two; this is known as dichotomous venation. The leaves are usually , but sometimes up to long. The old common name, maidenhair tree, derives from the leaves resembling pinnae of the maidenhair fern, Adiantum capillus-veneris. Ginkgos are prized for their autumn foliage, which is a deep saffron yellow. Leaves of long shoots are usually notched or lobed, but only from the outer surface, between the veins. They are borne both on the more rapidly growing branch tips, where they are alternate and spaced out, and also on the short, stubby spur shoots, where they are clustered at the tips. Leaves are green both on the top and bottom and have stomata on both sides. During autumn, the leaves turn a bright yellow and then fall, sometimes within a short space of time (one to fifteen days). Branches Ginkgo branches grow in length by growth of shoots with regularly spaced leaves, as seen on most trees. From the axils of these leaves, "spur shoots" (also known as short shoots) develop on second-year growth. Short shoots have short internodes (they may grow only one to two centimeters in several years) and their leaves are usually unlobed. They are short and knobby, and are arranged regularly on the branches except on first-year growth. Because of the short internodes, leaves appear to be clustered at the tips of short shoots, and reproductive structures are formed only on them (see pictures below – seeds and leaves are visible on short shoots). In ginkgos, as in other plants that possess them, short shoots allow the formation of new leaves in the older parts of the crown. After a number of years, a short shoot may change into a long (ordinary) shoot, or vice versa. Ginkgo prefers full sun and grows best in environments that are well-watered and well-drained. The species shows a preference for disturbed sites; in the "semiwild" stands at Tianmu Mountains, many specimens are found along stream banks, rocky slopes, and cliff edges. Accordingly, ginkgo retains a prodigious capacity for vegetative growth. It is capable of sprouting from embedded buds near the base of the trunk (lignotubers, or basal chichi) in response to disturbances, such as soil erosion. Old specimens are also capable of producing aerial roots on the undersides of large branches in response to disturbances such as crown damage; these roots can lead to successful clonal reproduction upon contacting the soil. These strategies are evidently important in the persistence of ginkgo; in a survey of the "semiwild" stands remaining in Tianmushan, 40% of the specimens surveyed were multi-stemmed, and few saplings were present. Reproduction Ginkgo biloba is dioecious, with separate sexes, some trees being female and others being male. Male plants produce small pollen cones with sporophylls, each bearing two microsporangia spirally arranged around a central axis. Sex conversion, wherein certain branches of a tree change sexes, has been observed. This phenomenon is difficult to research because of its rarity as well as the practice of grafting female branches onto otherwise male trees that was common in 19th century Europe. Female plants do not produce cones. Two ovules are formed at the end of a stalk, and after wind pollination, one or both develop into fruit-like structures containing seeds. The fruits are 1.5–2 cm long, with a soft, fleshy, yellow-brown outer layer (the sarcotesta) that is attractive in appearance, but contains butyric acid (also known as butanoic acid) and smells foul like rancid butter or vomit when fallen. Beneath the sarcotesta is the hard sclerotesta (the "shell" of the seed) and a papery endotesta, with the nucellus surrounding the female gametophyte at the center. The fertilization of ginkgo seeds occurs via motile sperm, as in cycads, ferns, mosses, and algae. The sperm are large (about 70–90 micrometres) and are similar to the sperm of cycads, which are slightly larger. Ginkgo sperm were first discovered by the Japanese botanist Sakugoro Hirase in 1896. The sperm have a complex multi-layered structure, which is a continuous belt of basal bodies that form the base of several thousand flagella which have a cilia-like motion. The flagella/cilia apparatus pulls the body of the sperm forwards. The sperm have only a tiny distance to travel to the archegonia, of which there are usually two or three. Two sperm are produced, one of which successfully fertilizes the ovule. Fertilization of ginkgo seeds occurs just before or after they fall in early autumn. Embryos may develop in the seeds before or after they drop from the tree. Genome Chinese scientists published a draft genome of Ginkgo biloba in 2016. The tree has a large genome of 10.6 billion DNA nucleobase "letters" (the human genome has three billion) and about 41,840 predicted genes which enable a considerable number of antibacterial and chemical defense mechanisms. 76.58% of the assembled sequence turned out to be repetitive sequences. In 2020, a study in China of ginkgo trees up to 667 years old showed little effects of aging, finding that the trees continued to grow with age and displayed no genetic evidence of senescence, and continued to make phytochemicals indefinitely. Phytochemicals Extracts of ginkgo leaves contain phenolic acids, proanthocyanidins, flavonoid glycosides, such as myricetin, kaempferol, isorhamnetin, and quercetin, and the terpene trilactones ginkgolides and bilobalides. The leaves also contain unique ginkgo biflavones, alkylphenols, and polyprenols. Taxonomy The older Chinese name for this plant is 銀果, meaning "silver fruit", pronounced yínguǒ in Mandarin or Ngan-gwo in Cantonese. The current commonly used names are 白果 (), meaning "white fruit", and (), meaning "silver apricot". The name 銀杏 was translated into Japanese as イチョウ () or ぎんなん () and into Korean as 은행 (). Carl Linnaeus described the species in 1771, the specific epithet biloba derived from the Latin bis, "twice" and loba, "lobed", referring to the shape of the leaves. Two names for the species recognise the botanist Richard Salisbury, a placement by Nelson as Pterophyllus salisburiensis and the earlier Salisburia adiantifolia proposed by James Edward Smith. The epithet of the latter may have been intended to denote a characteristic resembling Adiantum, the genus of maidenhair ferns. The scientific name Ginkgo is the result of a spelling error that occurred three centuries ago. Kanji typically have multiple pronunciations in Japanese, and the characters 銀杏 used for ginnan can also be pronounced ginkyō. Engelbert Kaempfer, the first Westerner to investigate the species in 1690, wrote down this pronunciation in the notes that he later used for the Amoenitates Exoticae (1712) with the "awkward" spelling "ginkgo". This appears to be a simple error of Kaempfer; taking his spelling of other Japanese words containing the syllable "kyō" into account, a more precise romanization following his writing habits would have been "ginkio" or "ginkjo". Linnaeus, who relied on Kaempfer when dealing with Japanese plants, adopted the spelling given in Kaempfer's "Flora Japonica" (Amoenitates Exoticae, p. 811). Kaempfer's drawing can be found in Hori's article. Classification The relationship of ginkgo to other plant groups remains uncertain. It has been placed loosely in the divisions Spermatophyta and Pinophyta, but no consensus has been reached. Since its seeds are not protected by an ovary wall, it can morphologically be considered a gymnosperm. The apricot-like structures produced by female ginkgo trees are technically not fruits, but are seeds that have a shell consisting of a soft and fleshy section (the sarcotesta), and a hard section (the sclerotesta). The sarcotesta has a strong smell that most people find unpleasant. The ginkgo is classified in its own division, the Ginkgophyta, comprising the single class Ginkgoopsida, order Ginkgoales, family Ginkgoaceae, genus Ginkgo and is the only extant species within this group. It is one of the best-known examples of a living fossil, because Ginkgoales other than G. biloba are not known from the fossil record after the Pliocene. Phylogeny Ginkgo biloba is a living fossil, with fossils recognisably related to modern ginkgo from the early Permian (Cisuralian), with likely oldest record being that of Trichopitys from the earliest Permian (Asselian) of France, over 290 million years old. The closest living relatives of the clade are the cycads, which share with the extant G. biloba the characteristic of motile sperm. Such plants with leaves that have more than four veins per segment have customarily been assigned to the taxon Ginkgo, while the taxon Baiera is used to classify those with fewer than four veins per segment. Sphenobaiera has been used for plants with a broadly wedge-shaped leaf that lacks a distinct leaf stem. Rise and decline Fossils attributable to the genus Ginkgo first appeared in the Middle Jurassic. The genus Ginkgo diversified and spread throughout Laurasia during the Jurassic and Early Cretaceous. The Ginkgophyta declined in diversity as the Cretaceous progressed, and by the Paleocene, Ginkgo adiantoides was the only Ginkgo species left in the Northern Hemisphere, while a markedly different (and poorly documented) form persisted in the Southern Hemisphere. Along with that of ferns, cycads, and cycadeoids, the species diversity in the genus Ginkgo drops through the Cretaceous, at the same time the flowering plants were on the rise; this supports the hypothesis that, over time, flowering plants with better adaptations to disturbance displaced Ginkgo and its associates. At the end of the Pliocene, Ginkgo fossils disappeared from the fossil record everywhere except in a small area of central China, where the modern species survived. Limited number of species It is doubtful whether the Northern Hemisphere fossil species of Ginkgo can be reliably distinguished. Given the slow pace of evolution and morphological similarity between members of the genus, there may have been only one or two species existing in the Northern Hemisphere through the entirety of the Cenozoic: present-day G. biloba (including G. adiantoides) and G. gardneri from the Paleocene of Scotland. At least morphologically, G. gardneri and the Southern Hemisphere species are the only known post-Jurassic taxa that can be unequivocally recognised. The remainder may have been ecotypes or subspecies. The implications would be that G. biloba had occurred over an extremely wide range, had remarkable genetic flexibility and, though evolving genetically, never showed much speciation. While it may seem improbable that a single species may exist as a contiguous entity for many millions of years, many of the ginkgo's life-history parameters fit: Extreme longevity; slow reproduction rate; (in Cenozoic and later times) a wide, apparently contiguous, but steadily contracting distribution; and (as far as can be demonstrated from the fossil record) extreme ecological conservatism (restriction to disturbed streamside environments). Adaptation to a single environment Given the slow rate of evolution of the genus, Ginkgo possibly represents a pre-angiosperm strategy for survival in disturbed streamside environments. Ginkgo evolved in an era before flowering plants, when ferns, cycads, and cycadeoids dominated disturbed streamside environments, forming low, open, shrubby canopies. Ginkgo large seeds and habit of "bolting" – growing to a height of 10 meters before elongating its side branches – may be adaptations to such an environment. Modern-day G. biloba grows best in environments that are well-watered and drained, and the extremely similar fossil Ginkgo favored similar environments: The sediment record at the majority of fossil Ginkgo localities indicates it grew primarily in disturbed environments, such as along streams. Ginkgo, therefore, presents an "ecological paradox" because while it possesses some favorable traits for living in disturbed environments (clonal reproduction) many of its other life-history traits are the opposite of those exhibited by modern plants that thrive in disturbed settings (slow growth, large seed size, late reproductive maturity). Etymology The genus name is regarded as a misspelling of the Japanese pronunciation () for the kanji meaning "silver apricot", which is found in Chinese herbology literature such as (Daily Use Materia Medica) (1329) and Compendium of Materia Medica published in 1578. Despite its spelling, which is due to a complicated etymology including a transcription error, "ginkgo" is usually pronounced , which has given rise to the common alternative spelling "gingko". The spelling pronunciation is also documented in some dictionaries. Engelbert Kaempfer first introduced the spelling ginkgo in his book of 1712. It is considered that he may have misspelled "Ginkjo" or "Ginkio" (both consistent with his treatment of Japanese in the same work) as "Ginkgo". This misspelling was included by Linnaeus in his book and has become the name of the tree's genus. The specific epithet is New Latin for "two-lobed". Distribution and habitat Although Ginkgo biloba and other species of the genus were once widespread throughout the world, its habitat had shrunk by two million years ago. For centuries, it was thought to be extinct in the wild, but is now a common tree cultivated throughout eastern China, Korea, and Japan. Many municipalities in Korea and Japan use Ginkgos as street trees, and Ginkgo leaves are the emblem of prominent educational institutions such as the University of Tokyo and Sungkyunkwan University in South Korea. Despite their widespread habitat, high genetic uniformity exists among ginkgo trees, with some Chinese scholars suggesting that ginkgo trees in these areas may have been planted and preserved by Chinese monks over about 1,000 years. A study demonstrates a greater genetic diversity in Southwestern China populations, supporting glacial refugia in mountains surrounding the eastern Tibetan Plateau, where several old-growth candidates for wild populations have been reported. Whether native ginkgo populations still exist has not been demonstrated unequivocally, but there is genetic evidence that these Southwestern populations may be wild, as well as evidence that the largest and oldest G. biloba trees may be older than surrounding human settlements. Where it occurs in the wild, Ginkgo is found infrequently in deciduous forests and valleys on acidic loess (i.e. fine, silty soil) with good drainage. The soil it inhabits is typically in the pH range of 5.0 to 5.5. Cultivation Ginkgo has long been cultivated in China. It is common in the southern third of the country. Some planted trees at temples are believed to be over 1,500 years old. The first record of Europeans encountering it is in 1690 in Japanese temple gardens, where the tree was seen by the German botanist Engelbert Kaempfer. Because of its status in Buddhism and Confucianism, the ginkgo has also been widely planted in Korea and in Japan since the 14th century; in both areas, some naturalization has occurred, with ginkgos seeding into natural forests. Ginkgo has been commonly cultivated in North America for over 200 years and in Europe for close to 300, but during that time, it has never become significantly naturalized. G. biloba is also commonly manually planted in cities across the United States and Europe. This species is highly tolerant to pollution and serves as a visually appealing, shade-providing tree in many cities and gardens. Many intentionally planted ginkgos are male cultivars grafted onto plants propagated from seed, because the male trees will not produce the malodorous seeds. The popular cultivar 'Autumn Gold' is a clone of a male plant. The disadvantage of male Ginkgo biloba trees is that they are highly allergenic. They have an OPALS (Ogren Plant Allergy Scale) rating of 7 (out of 10), whereas female trees, which can produce no pollen, have an OPALS allergy scale rating of 2. Female cultivars include 'Liberty Splendor', 'Santa Cruz', and 'Golden Girl', the latter so named because of the striking yellow color of its leaves in the fall; all female cultivars release zero pollen. Many cultivars are listed in the literature in the UK, of which the compact 'Troll' has gained the Royal Horticultural Society's Award of Garden Merit. Ginkgos adapt well to the urban environment, tolerating pollution and confined soil spaces. They rarely have disease problems, even in urban conditions, and are attacked by few insects. Ginkgos are popular subjects for growing as miniature landscapes known as penjing and bonsai; they can be kept artificially small and tended over centuries. The trees are easy to propagate from seed. Hiroshima Extreme examples of the ginkgo's tenacity may be seen in Hiroshima, Japan, where six trees growing between from the 1945 atom bomb explosion were among the few living organisms in the area to survive the blast. Although almost all other plants (and animals) in the area were killed, the ginkgos, though charred, survived and were soon healthy again, among other hibakujumoku (trees that survived the blast). The six trees are still alive: They are marked with signs at temple (planted in 1850), Shukkei-en (planted about 1740), Jōsei-ji (planted 1900), at the former site of Senda Elementary School near Miyukibashi, at the Myōjōin temple, and an Edo period-cutting at Anraku-ji temple. 1000-year-old ginkgo at Tsurugaoka Hachimangū At the Tsurugaoka Hachiman-gū's shrine in the city of Kamakura, Kanagawa Prefecture, Japan, an ancient ginkgo tree stands beside the stone entry staircase. According to legend, the tree has stood there since the founding of the shrine circa 1063. The tree is nicknamed kakure-ichō (hiding ginkgo), because of an Edo period legend in which shōgun Minamoto no Sanetomo was assassinated in 1219 by his nephew, Kugyō, who had hidden behind the tree to ambush the shōgun. Modern scholarship has established that ginkgos arrived from China in the 14th century, and a 1990 tree-ring measurement indicated the kakure-ichō's age to be about 500 years. On 10 March 2010, the tree blew down in a storm, but the stump has since sprouted vigorously. 1,400-year-old ginkgo tree at Gu Guanyin The grounds of the Buddhist temple at Gu Guanyin in the Zhongnan Mountains feature a ginkgo tree reputed to be 1,400 years old. The tree itself is a popular tourist attraction. Toxicity Since 2016, G. biloba extract is classified as a possible human carcinogen (group 2B) by the International Agency for Research on Cancer. When eaten in large quantities or over a long period, the seeds may cause poisoning by ginkgotoxin (4'-O-methylpyridoxine, MPN), as found in a few case reports. A heat-stable compound not destroyed by cooking, MPN may cause convulsions, which were alleviated by treatment with pyridoxine phosphate (vitamin B6), according to limited studies. Some people are sensitive to the chemicals in the sarcotesta, the outer fleshy coating. These people should handle the seeds with care when preparing the seeds for consumption, wearing disposable gloves. The symptoms are allergic contact dermatitis, or blisters similar to that caused by contact with poison ivy. Side effects of using ginkgo supplements may include increased risk of bleeding, gastrointestinal discomfort, nausea, vomiting, diarrhea, headaches, dizziness, heart palpitations, and restlessness. Although use of standardized Ginkgo biloba leaf extracts in moderate amounts appears to be safe, excessive use may have undesirable effects, especially in terms of drug interactions. The dosing of anticoagulants, such as warfarin or antiplatelet medication, may be adversely affected by using ginkgo supplements. According to a systemic review, the effects of ginkgo on pregnant women may include increased bleeding time, and there is inadequate information about safety during lactation. Ginkgo pollen may produce allergic reactions. Ginkgo biloba leaves and sarcotesta contain ginkgolic acids which are highly allergenic long-chain alkylphenols, such as bilobol or adipostatin A (bilobol is a substance related to anacardic acid from cashew nut shells and urushiols present in poison ivy and other Toxicodendron spp.) Individuals with a history of strong allergic reactions to poison ivy, mangoes, cashews and other alkylphenol-producing plants are more likely to experience an allergic reaction when consuming non-standardized ginkgo-containing preparations. The level of these allergens in standardized pharmaceutical preparations from Ginkgo biloba was restricted to 5 ppm by the Commission E of the former Federal German Health Authority. Overconsumption of seeds from Ginkgo biloba can deplete vitamin B6. Uses The wood of Ginkgo biloba is used to make furniture, chessboards, carving, and casks for making saké; the wood is fire-resistant and slow to decay. Culinary Despite the health risks in certain cases, the nut-like kernels of the seeds are esteemed in Asia, and are a traditional ingredient in Chinese food. Ginkgo nuts are used in congee, and are often served at special occasions such as weddings and the Chinese New Year (as part of the vegetarian dish called Buddha's delight). Japanese cooks add ginkgo seeds (called ginnan) to dishes such as chawanmushi, and cooked seeds are often eaten along with other dishes. Grilled ginkgo nuts with salt are also a popular item at izakayas as a healthy snack with beer and other Japanese food. In Korea, ginkgo nuts are stir-fried and eaten, or are used to garnish foods such as sinseonro. Medical research Although extracts of Ginkgo biloba leaf are often marketed as cognitive enhancers, there is no evidence for effects on memory or attention in healthy people. Systematic reviews have shown there is no evidence for effectiveness of ginkgo in treating high blood pressure, menopause-related cognitive decline, tinnitus, post-stroke recovery, or altitude sickness. There is weak preliminary evidence for ginkgo affecting dementia and tardive dyskinesia symptoms in people with schizophrenia. Traditional medicine Ginkgo has been used in traditional Chinese medicine since at least the 11th century CE. Ginkgo seeds, leaves, and nuts have traditionally been used to treat various ailments, such as dementia, asthma, bronchitis, and kidney and bladder disorders. However, there is no conclusive evidence that ginkgo is useful for any of these conditions. The European Medicines Agency Committee on Herbal Medicinal Products concluded that medicines containing ginkgo leaf can be used for treating mild age-related dementia and mild peripheral vascular disease in adults after serious conditions have been excluded by a physician. In culture The ginkgo leaf is the symbol of the Urasenke school of Japanese tea ceremony. The tree is the official tree of the Japanese capital of Tokyo, and the symbol of Tokyo is a ginkgo leaf. Since 1948, the badge of Tokyo University has been two ginkgo leaves (designed by Shoichi Hoshino), which became the university logo in 2004 with a redesign. The logo of Osaka University has been a simplified ginkgo leaf since 1991 when designer Ikko Tanaka created it for the university's sixtieth anniversary. In professional sumo, wrestlers ranked in the two highest divisions ( and ) wear an elaborate topknot called because it resembles the leaf of the ginkgo tree. Ginkgo is an official tree of Seoul since 1971, designated by the Seoul Metropolitan Government. Gallery	Biology and health sciences	Gymnosperms	null
67970	https://en.wikipedia.org/wiki/Lycophyte	Lycophyte	The lycophytes, when broadly circumscribed, are a group of vascular plants that include the clubmosses. They are sometimes placed in a division Lycopodiophyta or Lycophyta or in a subdivision Lycopodiophytina. They are one of the oldest lineages of extant (living) vascular plants; the group contains extinct plants that have been dated from the Silurian (ca. 425 million years ago). Lycophytes were some of the dominating plant species of the Carboniferous period, and included the tree-like Lepidodendrales, some of which grew over in height, although extant lycophytes are relatively small plants. The scientific names and the informal English names used for this group of plants are ambiguous. For example, "Lycopodiophyta" and the shorter "Lycophyta" as well as the informal "lycophyte" may be used to include the extinct zosterophylls or to exclude them. Description Lycophytes reproduce by spores and have alternation of generations in which (like other vascular plants) the sporophyte generation is dominant. Some lycophytes are homosporous while others are heterosporous. When broadly circumscribed, the lycophytes represent a line of evolution distinct from that leading to all other vascular plants, the euphyllophytes, such as ferns, gymnosperms and flowering plants. They are defined by two synapomorphies: lateral rather than terminal sporangia (often kidney-shaped or reniform), and exarch protosteles, in which the protoxylem is outside the metaxylem rather than vice versa. The extinct zosterophylls have at most only flap-like extensions of the stem ("enations") rather than leaves, whereas extant lycophyte species have microphylls, leaves that have only a single vascular trace (vein), rather than the much more complex megaphylls of other vascular plants. The extinct genus Asteroxylon represents a transition between these two groups: it has a vascular trace leaving the central protostele, but this extends only to the base of the enation. See . Zosterophylls and extant lycophytes are all relatively small plants, but some extinct species, such as the Lepidodendrales, were tree-like, and formed extensive forests that dominated the landscape and contributed to the formation of coal. Taxonomy Classification In the broadest circumscription of the lycophytes, the group includes the extinct zosterophylls as well as the extant (living) lycophytes and their closest extinct relatives. The names and ranks used for this group vary considerably. Some sources use the names "Lycopodiophyta" or the shorter "Lycophyta" to include zosterophylls as well as extant lycophytes and their closest extinct relatives, while others use these names to exclude zosterophylls. The name "Lycopodiophytina" has also been used in the inclusive sense. English names, such as "lycophyte", "lycopodiophyte" or "lycopod", are similarly ambiguous, and may refer to the broadly defined group or only to the extant lycophytes and their closest extinct relatives. The consensus classification produced by the Pteridophyte Phylogeny Group classification in 2016 (PPG I) places all extant (living) lycophytes in the class Lycopodiopsida. There are around 1,290 to 1,340 such species. For more information on the classification of extant lycophytes, see . Phylogeny A major cladistic study of land plants was published in 1997 by Kenrick and Crane. In 2004, Crane et al. published some simplified cladograms, based on a number of figures in Kenrick and Crane (1997). Their cladogram for the lycophytes is reproduced below (with some branches collapsed into 'basal groups' to reduce the size of the diagram). In this view, the "zosterophylls" comprise a paraphyletic group, ranging from forms like Hicklingia, which had bare stems, to forms like Sawdonia and Nothia, whose stems are covered with unvascularized spines or enations. The genus Renalia illustrates the problems in classifying early land plants. It has characteristics both of the non-lycophyte rhyniophytes – terminal rather than lateral sporangia – and of the zosterophylls – kidney-shaped sporangia opening along the distal margin. A rather different view is presented in a 2013 analysis by Hao and Xue. Their preferred cladogram shows the zosterophylls and associated genera basal to both the lycopodiopsids and the euphyllophytes, so that there is no clade corresponding to the broadly defined group of lycophytes used by other authors. Some extinct orders of lycophytes fall into the same group as the extant orders. Different sources use varying numbers and names of the extinct orders. The following phylogram shows a likely relationship between some of the proposed Lycopodiopsida orders. Evolution of microphylls Within the broadly defined lycophyte group, species placed in the class Lycopodiopsida are distinguished from species placed in the Zosterophyllopsida by the possession of microphylls. Some zosterophylls, such as the Devonian Zosterophyllum myretonianum, had smooth stems (axes). Others, such as Sawdonia ornata, had flap-like extensions on the stems ("enations"), but without any vascular tissue. Asteroxylon, identified as an early lycopodiopsid, had vascular traces that extended to the base of the enations. Species in the genus Leclercqia had fully vascularized microphylls. These are considered to be stages in the evolution of microphylls. Gallery	Biology and health sciences	Pteridophytes	null
68005	https://en.wikipedia.org/wiki/Loon	Loon	Loons (North American English) or divers (British / Irish English) are a group of aquatic birds found in much of North America and northern Eurasia. All living species of loons are members of the genus Gavia, family Gaviidae and order Gaviiformes. Description Loons, which are the size of large ducks or small geese, resemble these birds in shape when swimming. Like ducks and geese, but unlike coots (which are Rallidae) and grebes (Podicipedidae), the loon's toes are connected by webbing. The loons may be confused with the cormorants (Phalacrocoracidae), but can be distinguished from them by their distinct call. Cormorants are not-too-distant relatives of loons, and like them are heavy-set birds whose bellies, unlike those of ducks and geese, are submerged when swimming. Loons in flight resemble plump geese with seagulls' wings that are relatively small in proportion to their bulky bodies. The bird points its head slightly upwards while swimming, but less so than cormorants. In flight, the head droops more than in similar aquatic birds. Male and female loons have identical plumage, which is largely patterned black-and-white in summer, with grey on the head and neck in some species. All have a white belly. This resembles many sea-ducks (Merginae) – notably the smaller goldeneyes (Bucephala) – but is distinct from most cormorants, which rarely have white feathers, and if so, usually as large rounded patches rather than delicate patterns. All species of loons have a spear-shaped bill. Males are larger on average, but relative size is only apparent when the male and female are together. In winter, plumage is dark grey above, with some indistinct lighter mottling on the wings, and a white chin, throat and underside. The specific species can then be distinguished by certain features, such as the size and colour of the head, neck, back and bill. But reliable identification of loons in winter is often difficult even for experts – particularly as the smaller immature birds look similar to winter-plumage adults, making size an unreliable means of identification. Gaviiformes are among the few groups of birds in which the young moult into a second coat of down feathers after shedding the first one, rather than growing juvenile feathers with downy tips that wear off, as is typical in many birds. This trait is also found in tubenoses (Procellariiformes) and penguins (Sphenisciformes), both relatives of the loons. Behaviour and ecology Loons are excellent swimmers, using their feet to propel themselves above and under water. However, since their feet are located far back on the body, loons have difficulty walking on land, though they can effectively run short distances to reach water when frightened. Thus, loons avoid coming to land, except for mating and nesting. Loons fly strongly, though they have high wing loading (mass to wing area ratio), which complicates takeoff. Indeed, most species must run upwind across the water's surface with wings flapping to generate sufficient lift to take flight. Only the red-throated loon (G. stellata) can take off from land. Once airborne, loons are capable of long flights during migration. Scientists from the U.S. Geological Survey, who have implanted satellite transmitters in some individuals, have recorded daily flights of up to 1078 km in a 24-hour period, which probably resulted from single movements. North European loons migrate primarily via the South Baltic and directly over land to the Black Sea or Mediterranean. Loons can live as long as 30 years and can hold their breath for as long as 90 seconds while underwater. Loons are migratory birds, and in the winter months they move from their northern freshwater lake nesting habitats to southern marine coastlines. They are well-adapted to this change in salinity, however, because they have special salt glands located directly above their eyes. These glands filter out salts in their blood and flush this salty solution out through their nasal passages, which allows them to immediately consume fish from oceans and drink saltwater after their long migration. Diet and feeding Loons find their prey by sight. They eat mainly fish, supplemented with amphibians, crustaceans and similar mid-sized aquatic fauna. Specifically, they have been noted to feed on crayfish, frogs, snails, salamanders and leeches. They prefer clear lakes because they can more easily see their prey through the water. The loon uses its pointy bill to stab or grasp prey. They eat vertebrate prey headfirst to facilitate swallowing, and swallow all their prey whole. To help digestion, loons swallow small pebbles from the bottoms of lakes. Similar to grit eaten by chickens, these gastroliths may assist the loon's gizzard in crushing the hard parts of the loon's food such as the exoskeletons of crustaceans and the bones of frogs and salamanders. The gastroliths may also be involved in stomach cleaning as an aid to regurgitation of indigestible food parts. Loons may inadvertently ingest small lead pellets, released by anglers and hunters, that will contribute to lead poisoning and the loon's eventual death. Jurisdictions that have banned the use of lead shot and sinkers include but are not limited to Maine, New Hampshire, Vermont, Michigan, some areas of Massachusetts, Yellowstone National Park, Canada, Great Britain, and Denmark. Reproduction Loons nest during the summer on freshwater lakes and/or large ponds. Smaller bodies of water (up to 0.5 km2) will usually only have one pair. Larger lakes may have more than one pair, with each pair occupying a bay or section of the lake. The red-throated loon, however, may nest colonially, several pairs close together, in small Arctic tarns and feed at sea or in larger lakes, ferrying the food in for the young. Loons mate on land, often on the future nest site, and build their nests close to the water, preferring sites that are completely surrounded by water such as islands or emergent vegetation. Loons use a variety of materials to build their nests including aquatic vegetation, pine needles, leaves, grass, moss and mud. Sometimes, nest material is almost lacking. Both male and female build the nest and incubate jointly for 28 days. If the eggs are lost, the pair may re-nest, usually in a different location. Since the nest is very close to the water, rising water may induce the birds to slowly move the nest upwards, over a metre. Despite the roughly equal participation of the sexes in nest building and incubation, analysis has shown clearly that males alone select the location of the nest. This pattern has the important consequence that male loons, but not females, establish significant site-familiarity with their territories that allows them to produce more chicks there over time. Sex-biased site-familiarity might explain, in part, why resident males fight so hard to defend their territories. Most clutches consist of two eggs, which are laid in May or June, depending upon latitude. Loon chicks are precocial, able to swim and dive right away, but will often ride on their parents' back during their first two weeks to rest, conserve heat, and avoid predators. Chicks are fed mainly by their parents for about six weeks but gradually begin to feed themselves over time. By 11 or 12 weeks, chicks gather almost all of their own food and have begun to fly. In 2019, a necropsy of a bald eagle found floating on a Maine lake (beside the floating body of a loon chick) found that the eagle had been stabbed through the heart by an adult loon's beak. Biologists, especially from Chapman University, have extensively studied the mating behaviour of the common loon (G. immer). Contrary to popular belief, pairs seldom mate for life. Indeed, a typical adult loon is likely to have several mates during its lifetime because of territorial takeover. Each breeding pair must frequently defend its territory against "floaters" (territory-less adults) trying to evict at least one owner and seize the breeding site. Territories that have produced chicks in the past year are especially prone to takeovers, because nonbreeding loons use chicks as cues to indicate high-quality territories. One-third of all territorial evictions among males result in the death of the owner; in contrast, female loons usually survive. Birds that are displaced from a territory but survive usually try to re-mate and (re)claim a breeding territory later in life. In 2020, a loon hatched for the first time in over a century in Southeastern Massachusetts at Fall River, the Massachusetts Division of Fisheries and Wildlife and Biodiversity Research Institute. The chicks were relocated in 2015 with the hopes of re-establishing breeding and nesting patterns. Etymology and taxonomy The European Anglophone name "diver" comes from the bird's habit of catching fish by swimming calmly along the surface and then abruptly plunging into the water. The North American name "loon" likely comes from either the Old English word lumme, meaning lummox or awkward person, or the Scandinavian word lum meaning lame or clumsy. Either way, the name refers to the loon's poor ability to walk on land. Another possible derivation is from the Norwegian word lom for these birds, which comes from Old Norse lómr, possibly cognate with English "lament", referring to the characteristic plaintive sound of the loon. The scientific name Gavia refers to seabirds in general. The scientific name Gavia was the Latin term for the smew (Mergellus albellus). This small sea-duck is quite unrelated to loons and just happens to be another black-and-white seabird which swims and dives for fish. It is not likely that the ancient Romans had much knowledge of loons, as these are limited to more northern latitudes and since the end of the last glacial period seem to have occurred only as rare winter migrants in the Mediterranean region. The term gavia was transferred from the ducks to the loons only in the 18th century. Earlier naturalists referred to the loons as mergus (the Latin term for diving seabirds of all sorts) or colymbus, which became the genus name used in the first modern scientific description of a Gavia species (by Carl Linnaeus) in 1758. Unfortunately, confusion about whether Linnaeus' "wastebin genus" Colymbus referred to loons or grebes abounded. North American ornithologists used the genus name to refer to grebes, while Europeans used it for loons, following Nicholas Aylward Vigors and Richard Bowdler Sharpe. The International Commission on Zoological Nomenclature tried to settle this issue in 1956 by declaring Colymbus a suppressed name unfit for further use and establishing Gavia, created by Johann Reinhold Forster in 1788, as the valid genus name for the loons. However, the situation was not completely resolved even then, and the following year the ICZN had to act again to prevent Louis Pierre Vieillot's 1818 almost-forgotten family name Urinatoridae from overruling the much younger Gaviidae. Some eminent ornithologists such as Pierce Brodkorb tried to keep the debate alive, but the ICZN's solution has been satisfactory. Systematics and evolution All living species are placed in the genus Gavia. It has been suggested that the genus Gavia originated in Europe during the Paleogene. The earliest species, G. egeriana, was found in early Miocene deposits in Dolnice in the Czech Republic. During the remainder of the Miocene, Gavia managed to disperse into North America via the Atlantic coastlines, eventually making their way to the continent's Pacific coastlines by the Late Miocene. Study of the interrelationships of the extant species has found that the red-throated loons are the most basal of the five species. Fossil record Nearly ten prehistoric species have been named to date in the genus Gavia, and about as many undescribed ones await further study. The genus is known from the Early Miocene onwards, and the oldest members are rather small (some are smaller than the red-throated loon). Throughout the late Neogene, the genus by and large follows Cope's Rule (that population lineages tend to increase in body size over evolutionary time). List of fossil Gavia species †G. brodkorbi Howard, 1978 (Late Miocene of Orange County, United States) †G. concinna Wetmore, 1940 (Late Miocene/Early Pliocene of west and east United States) †G. egeriana Švec, 1982 (Early Miocene of Czechoslovakia ?and Cheswold, Delaware, United States –? Yorktown Early Pliocene of Lee Creek Mine, North Carolina, United States) †G. fortis Olson & Rasmussen, 2001 (Yorktown Early Pliocene of Lee Creek Mine, North Carolina, United States) †G. howardae Brodkorb, 1953 (San Diego Formation, California and Yorktown Formation, North Carolina †G. moldavica Kessler, 1984 (Late Miocene of Chişinău, Moldova) †G. palaeodytes Wetmore, 1943 (Bone Valley Early/Middle Pliocene of Pierce, Florida, United States) †G. paradoxa Umanska, 1981 (Late Miocene of Čebotarevka, Ukraine) †G. schultzi Mlíkovský, 1998 (Middle Miocene of Sankt Margarethen, Austria) List of fossil Gavia specimens Gavia sp. (Early-Middle Miocene of eastern United States) Gavia sp. (Calvert Middle Miocene ?or Pleistocene of Maryland, United States) – same as Gavia cf. immer below? Gavia spp. (Middle Miocene of Steinheim, Germany) – three species Gavia sp. (Early Pliocene of Empoli, Italy) Gavia sp. (Early Pliocene of Kerč Peninsula, Ukraine) Gavia cf. concinna (San Diego Middle/Late Pliocene of San Diego, California, United States) – two species? Gavia sp. (Early Pleistocene of Kairy, Ukraine) Gavia cf. immer (Pleistocene of California and Florida, United States) – possibly a G. immer paleosubspecies "Gavia" portisi from the Late Pliocene of Orciano Pisano, Italy, is known from a cervical vertebra that may or may not have been from a loon. If so, it was from a bird slightly smaller than the common loon. Older authors were quite sure the bone was indeed from a Gavia and even considered G. concinna a possibly junior synonym of it. This is now regarded as rather unlikely due to the quite distinct range and age. The Early Pliocene Gavia skull from Empoli (Italy) was referred to G. concinna, and thus could conceivably have been of "G." portisi if that was indeed a loon. The holotype vertebra may now be lost, which would make "G." portisi a nomen dubium. In popular culture Various Indigenous myths from the California region have a recurring figure, Loon or Loon Woman, based on the common loon. The Tlingit of Alaska believe that loon cries forecast rain. Gaǥeit is named after the common loon (kagit). The common loon is the provincial bird of Ontario and is depicted on the Canadian one-dollar coin, which has come to be known affectionately as the "loonie". The common loon is the official state bird of Minnesota. Mercer, Wisconsin, promotes itself as the "Loon Capital of the World". Henry David Thoreau describes a playful and inspiring acquaintance with a loon on Walden Pond in his book Walden. The Great Lakes Loons are a minor-league professional baseball team based in Midland, Michigan, United States. The primary mascot is Lou E. Loon. The Warner Bros./Amblin cartoon Tiny Toon Adventures features Shirley the Loon, who speaks with a thick Valley girl accent and is obsessed with superficial New Age paraphernalia. She is voiced by Saturday Night Live cast member Gail Matthius. The Major League Soccer club Minnesota United FC use a loon in the club's crest, and as a nickname for the team. Thanks to its inclusion as a preset in the E-mu Emulator, a specific sample of a Canadian loon, notably heard in "Sueño Latino" (1989) and in 808 State’s "Pacific State" (1989), has become a recurring motif in electronic-based popular music.	Biology and health sciences	Gaviiformes	null
68073	https://en.wikipedia.org/wiki/Welwitschia	Welwitschia	Welwitschia is a monotypic genus (that is, a genus that contains a single recognised species) of gymnosperm, the sole described species being the distinctive Welwitschia mirabilis, endemic to the Namib desert within Namibia and Angola. Welwitschia is the only extant genus of the family Welwitschiaceae and order Welwitschiales in the division Gnetophyta, and is one of three extant genera in Gnetophyta, alongside Gnetum and Ephedra. Informal sources commonly refer to the plant as a "living fossil". Naming Welwitschia is named after Austrian botanist and doctor Friedrich Welwitsch, who described the plant in Angola in 1859. Welwitsch was so overwhelmed by the plant that he "could do nothing but kneel down [...] and gaze at it, half in fear lest a touch should prove it a figment of the imagination." Joseph Dalton Hooker of the Linnean Society of London described the species, using Welwitsch's description and collected material along with material from artist Thomas Baines who had independently recorded the plant in Namibia. Welwitsch proposed calling the genus Tumboa after what he believed to be the local name, tumbo. Hooker asked Welwitsch for permission to name the genus Welwitschia instead. Welwitsch concurred and supplied some well-preserved material from which Hooker was able to make substantial progress in determining its botanical affinities. The taxonomy of Welwitschia subsequently changed intermittently with the development of new classification systems (see Flowering plants: History of classification); however, its current taxonomic status is essentially the same as Hooker's placement. Most botanists have treated Welwitschia as a distinct monotypic genus in a monotypic family or even order. Most recent systems place Welwitschia mirabilis in its own family Welwitschiaceae in the gymnosperm order Gnetales, although other extinct species have been placed in this family. The plant is commonly known simply as welwitschia in English, but the name tree tumbo is also used. It is called or in Nama, ('two leaves; can't die') in Afrikaans, in Damara, and in Herero. Biology After germination, the seedling produces two cotyledons which grow to in length, and have reticulate venation. Subsequently, two foliage leaves are produced at the edge of a woody bilobed crown. The permanent leaves are opposite (at right angles to the cotyledons), amphistomatic (producing stomata on both sides of the leaf), parallel-veined and ribbon-shaped. Shortly after the appearance of the foliage leaves, the apical meristem dies and meristematic activity is transferred to the periphery of the crown. The two (rarely three) foliage leaves grow continuously from a basal meristem around the circumference of the trunk, reaching lengths up to . The tips of the leaves split and fray into several well-separated strap-shaped sections by the distortions of the woody portions surrounding the apical slit, and also by wind and adventitious external injuries. The largest specimens (such as the "Husab Giant" which is five meters in circumference (about five feet in diameter)) may be no more than tall above ground, but the circumference of the leaves in contact with the sand may exceed . Welwitschia has an elongated shallow root system consisting of "a tapering taproot with one or more non-tapering extensions, some pronounced lateral roots, and a network of delicate spongy roots" and a woody fibrous unbranched main stem. The roots extend to a depth roughly equal to the span of the living leaves from tip to tip. The main stem consists of an unbranched woody crown roughly shaped like an inverted cone. The only branching in the shoot system occurs in the reproductive branches, which bear strobili. The species is dioecious, with separate male and female plants. Fertilization is carried out by insects including flies and true bugs. The most common of the true bugs attending Welwitschia is a member of the family Pyrrhocoridae, Probergrothius angolensis, but a hypothesized role in pollination has so far not been demonstrated. Infrequently, wasps and bees also play a role as pollinators of Welwitschia. At least some of the pollinators are attracted by "nectar" produced on both male and female strobili. Welwitschia has been classified as a CAM plant (crassulacean acid metabolism) after reconciliation of some initially contradictory and confusing data. There are however some very puzzling aspects to the matter; for example, the employment of the CAM metabolism is very slight, which was part of the reason that it took so long to establish its presence at all; it is not understood why this should be. The age of individual plants is difficult to assess, but many plants may be over 1,000 years old. Some individuals may be more than 2,000 years old. As the species does not produce yearly rings, plant age is determined by radiocarbon dating. However, other reports suggest that the plant does produce a kind of yearly ring. The "trunk" continues to expand with age. The largest known is in diameter ( in circumference). Because Welwitschia only produces a single pair of foliage leaves, the plant was thought by some to be neotenic, consisting essentially of a "giant seedling." However, research showed that its anatomy is not consistent with the giant seedling idea. Instead, the plant is more accurately thought to achieve its unusual morphology as a result of having "lost its head" (apical meristem) at an early stage. Genetics In July 2021, the genome of Welwitschia was 98% sequenced, totaling 6.8 Gb on 21 chromosomes. There is evidence of a whole genome duplication followed by extensive reshuffling, probably caused by extreme stress due to a time of increased aridity and prolonged drought some 86 million years ago. As a result of this duplication, the genome contains more "junk" self-replicating DNA sequences. This increase in retrotransposon activity was counteracted with a silencing DNA methylation process allowing to lower the metabolic cost of such a large genetic material and improve resilience. Distribution and habitat W. mirabilis is endemic to the desert bordering the Angolan and Namibian coast, between and inland, and from 14.12°S, near the Bentiaba River in Angola, to 23.64°S, near the Kuiseb River in Namibia, a distance of . The area is arid; the coast is recorded as having almost zero rainfall, while less than of rain falls annually below the escarpment in the wet season from February to April. Populations tend to occur in ephemeral watercourses, indicating a dependence on groundwater in addition to precipitation from fog. Cultivation Welwitschia mirabilis grows readily from seed, which may be bought from specialty seed dealers. The seeds have been shown to display orthodox seed behavior, which in general means that they may be stored for long periods at suitably low humidity and temperature. Welwitschia seeds naturally develop suitably low water concentrations as they ripen. Removal of the outer seed coverings enhances germination performance, which suggests that the seeds may display non-deep physiological dormancy. On planting the seed it is necessary to keep it moist but not immersed in water for the first two weeks of cultivation; it has been suggested that soaking the seeds in water before planting interferes with germination. Seeds collected from the wild often are heavily contaminated with spores of the fungus Aspergillus niger var. phoenicis, which causes them to rot shortly after they germinate. The fungal inoculum infects the growing cones of W. mirabilis early during their development, and a sharp increase in infection occurs when the pollination drops appear; through those drops the fungal spores may gain access to the interior of the developing seed. Seeds in the wild may therefore be obliterated through fungal action even before they are fully developed. Seeds from botanical gardens or other cultivated sources are much cleaner and less likely to rot. The fungicide tebuconazole may be useful in controlling limited A. niger seed infection. As food Indigenous people eat the cone of this plant by eating it raw or baking it in hot ashes. One of its names, onyanga, translates to 'onion of the desert'. Conservation The population of Welwitschia mirabilis in the wild is reasonably satisfactory at present. The international trade in the plant is controlled under the Convention on International Trade in Endangered Species of Wild Fauna and Flora (CITES). Plants in Angola are better protected than those in Namibia, because the relatively high concentration of land mines in Angola keep collectors away. Although Welwitschia mirabilis is not at present immediately threatened, there being abundant populations over a large area, its status is far from secure; its recruitment and growth rates are low, and its range, though wide, covers only a single compact, ecologically limited and vulnerable area. The remarkable longevity of Welwitschia favours its survival of temporary periods adverse to reproduction, but it offers no protection against circumstances of direct threat, such as overgrazing and disease. Fungal infection of female cones severely reduces seed viability, reducing already inherently low recruitment. Other threats include injury from off-road vehicles, collection of wild plants and overgrazing by zebras, rhinos, and domestic animals. Heraldry The plant figures in the compartment of the national coat of arms of Namibia. Gallery	Biology and health sciences	Gymnosperms (except conifers)	Plants
68085	https://en.wikipedia.org/wiki/Conifer	Conifer	Conifers are a group of cone-bearing seed plants, a subset of gymnosperms. Scientifically, they make up the division Pinophyta (), also known as Coniferophyta () or Coniferae. The division contains a single extant class, Pinopsida. All extant conifers are perennial woody plants with secondary growth. The great majority are trees, though a few are shrubs. Examples include cedars, Douglas-firs, cypresses, firs, junipers, kauri, larches, pines, hemlocks, redwoods, spruces, and yews. As of 2002, Pinophyta contained seven families, 60 to 65 genera, and more than 600 living species. Although the total number of species is relatively small, conifers are ecologically important. They are the dominant plants over large areas of land, most notably the taiga of the Northern Hemisphere, but also in similar cool climates in mountains further south. Boreal conifers have many wintertime adaptations. The narrow conical shape of northern conifers, and their downward-drooping limbs, help them shed snow. Many of them seasonally alter their biochemistry to make them more resistant to freezing. While tropical rainforests have more biodiversity and turnover, the immense conifer forests of the world represent the largest terrestrial carbon sink. Conifers are of great economic value for softwood lumber and paper production. Names and taxonomy Conifer is a Latin word, a compound of conus (cone) and ferre (to bear), meaning "the one that bears (a) cone(s)". The division name Pinophyta conforms to the rules of the International Code of Nomenclature for algae, fungi, and plants (ICN), which state (Article 16.1) that the names of higher taxa in plants (above the rank of family) are either formed from the name of an included family (usually the most common and/or representative), in this case Pinaceae (the pine family), or are descriptive. A descriptive name in widespread use for the conifers (at whatever rank is chosen) is Coniferae (Art 16 Ex 2). According to the ICN, it is possible to use a name formed by replacing the termination -aceae in the name of an included family, in this case preferably Pinaceae, by the appropriate termination, in the case of this division -ophyta. Alternatively, "descriptive botanical names" may also be used at any rank above family. Both are allowed. This means that if conifers are considered a division, they may be called Pinophyta or Coniferae. As a class, they may be called Pinopsida or Coniferae. As an order they may be called Pinales or Coniferae or Coniferales. Conifers are the largest and economically most important component group of gymnosperms, but nevertheless they comprise only one of the four groups. The division Pinophyta consists of just one class, Pinopsida, which includes both living and fossil taxa. Subdivision of the living conifers into two or more orders has been proposed from time to time. The most commonly seen in the past was a split into two orders, Taxales (Taxaceae only) and Pinales (the rest), but recent research into DNA sequences suggests that this interpretation leaves the Pinales without Taxales as paraphyletic, and the latter order is no longer considered distinct. A more accurate subdivision would be to split the class into three orders, Pinales containing only Pinaceae, Araucariales containing Araucariaceae and Podocarpaceae, and Cupressales containing the remaining families (including Taxaceae), but there has not been any significant support for such a split, with the majority of opinion preferring retention of all the families within a single order Pinales, despite their antiquity and diverse morphology. There were seven families of conifers , with 65–70 genera and over 600 living species (). The seven most distinct families are linked in the box above right and phylogenetic diagram left. In other interpretations, the Cephalotaxaceae may be better included within the Taxaceae, and some authors additionally recognize Phyllocladaceae as distinct from Podocarpaceae (in which it is included here). The family Taxodiaceae is here included in the family Cupressaceae, but was widely recognized in the past and can still be found in many field guides. A new classification and linear sequence based on molecular data can be found in an article by Christenhusz et al. The conifers are an ancient group, with a fossil record extending back about 300 million years to the Paleozoic in the late Carboniferous period; even many of the modern genera are recognizable from fossils 60–120 million years old. Other classes and orders, now long extinct, also occur as fossils, particularly from the late Paleozoic and Mesozoic eras. Fossil conifers included many diverse forms, the most dramatically distinct from modern conifers being some herbaceous conifers with no woody stems. Major fossil orders of conifers or conifer-like plants include the Cordaitales, Vojnovskyales, Voltziales and perhaps also the Czekanowskiales (possibly more closely related to the Ginkgophyta). Multiple studies also indicate that the Gnetophyta belong within the conifers despite their distinct appearances, either placing them as a sister group to Pinales (the 'gnepine' hypothesis) or as being more derived than Pinales but sister to the rest of the group. Most recent studies favor the 'gnepine' hypothesis. Phylogeny The earliest conifers appear in the fossil record during the Late Carboniferous (Pennsylvanian), over 300 million years ago. Conifers are thought to be most closely related to the Cordaitales, a group of extinct Carboniferous-Permian trees and clambering plants whose reproductive structures had some similarities to those of conifers. The most primitive conifers belong to the paraphyletic assemblage of "walchian conifers", which were small trees, and probably originated in dry upland habitats. The range of conifers expanded during the Early Permian (Cisuralian) to lowlands due to increasing aridity. Walchian conifers were gradually replaced by more advanced voltzialean or "transition" conifers. Conifers were largely unaffected by the Permian–Triassic extinction event, and were dominant land plants of the Mesozoic era. Modern groups of conifers emerged from the Voltziales during the Late Permian through Jurassic. Conifers underwent a major decline in the Late Cretaceous corresponding to the explosive adaptive radiation of flowering plants. Description All living conifers are woody plants, and most are trees, the majority having a monopodial growth form (a single, straight trunk with side branches) with strong apical dominance. Many conifers have distinctly scented resin, secreted to protect the tree against insect infestation and fungal infection of wounds. Fossilized resin hardens into amber, which has been commercially exploited historically (for example, in New Zealand's 19th-century kauri gum industry). The size of mature conifers varies from less than one metre to over 100 metres in height. The world's tallest, thickest, largest, and oldest living trees are all conifers. The tallest is a coast redwood (Sequoia sempervirens), with a height of 115.55 metres (although one mountain ash, Eucalyptus regnans, allegedly grew to a height of 140 metres, the tallest living angiosperms are significantly smaller at around 100 metres.) The thickest (that is, the tree with the greatest trunk diameter) is a Montezuma cypress (Taxodium mucronatum), 11.42 metres in diameter. The largest tree by three-dimensional volume is a giant sequoia (Sequoiadendron giganteum), with a volume 1486.9 cubic metres. The smallest is the pygmy pine (Lepidothamnus laxifolius) of New Zealand, which is seldom taller than 30 cm when mature. The oldest non-clonal living tree is a Great Basin bristlecone pine (Pinus longaeva), 4,700 years old. Foliage Since most conifers are evergreens, the leaves of many conifers are long, thin and have a needle-like appearance, but others, including most of the Cupressaceae and some of the Podocarpaceae, have flat, triangular scale-like leaves. Some, notably Agathis in Araucariaceae and Nageia in Podocarpaceae, have broad, flat strap-shaped leaves. Others such as Araucaria columnaris have leaves that are awl-shaped. In the majority of conifers, the leaves are arranged spirally, the exceptions being most of Cupressaceae and one genus in Podocarpaceae, where they are arranged in decussate opposite pairs or whorls of 3 (−4). In many species with spirally arranged leaves, such as Abies grandis (pictured), the leaf bases are twisted to present the leaves in a very flat plane for maximum light capture. Leaf size varies from 2 mm in many scale-leaved species, up to 400 mm long in the needles of some pines (e.g. Apache pine, Pinus engelmannii). The stomata are in lines or patches on the leaves and can be closed when it is very dry or cold. The leaves are often dark green in colour, which may help absorb a maximum of energy from weak sunshine at high latitudes or under forest canopy shade. Conifers from hotter areas with high sunlight levels (e.g. Turkish pine Pinus brutia) often have yellower-green leaves, while others (e.g. blue spruce, Picea pungens) may develop blue or silvery leaves to reflect ultraviolet light. In the great majority of genera the leaves are evergreen, usually remaining on the plant for several (2–40) years before falling, but five genera (Larix, Pseudolarix, Glyptostrobus, Metasequoia and Taxodium) are deciduous, shedding their leaves in autumn. The seedlings of many conifers, including most of the Cupressaceae, and Pinus in Pinaceae, have a distinct juvenile foliage period where the leaves are different, often markedly so, from the typical adult leaves. Tree ring structure Tree rings are records of the influence of environmental conditions, their anatomical characteristics record growth rate changes produced by these changing conditions. The microscopic structure of conifer wood consists of two types of cells: parenchyma, which have an oval or polyhedral shape with approximately identical dimensions in three directions, and strongly elongated tracheids. Tracheids make up more than 90% of timber volume. The tracheids of earlywood formed at the beginning of a growing season have large radial sizes and smaller, thinner cell walls. Then, the first tracheids of the transition zone are formed, where the radial size of cells and the thickness of their cell walls changes considerably. Finally, latewood tracheids are formed, with small radial sizes and greater cell wall thickness. This is the basic pattern of the internal cell structure of conifer tree rings. Reproduction Most conifers are monoecious, but some are subdioecious or dioecious; all are wind-pollinated. Conifer seeds develop inside a protective cone called a strobilus. The cones take from four months to three years to reach maturity, and vary in size from long. In Pinaceae, Araucariaceae, Sciadopityaceae and most Cupressaceae, the cones are woody, and when mature the scales usually spread open allowing the seeds to fall out and be dispersed by the wind. In some (e.g. firs and cedars), the cones disintegrate to release the seeds, and in others (e.g. the pines that produce pine nuts) the nut-like seeds are dispersed by birds (mainly nutcrackers, and jays), which break up the specially adapted softer cones. Ripe cones may remain on the plant for a varied amount of time before falling to the ground; in some fire-adapted pines, the seeds may be stored in closed cones for up to 60–80 years, being released only when a fire kills the parent tree. In the families Podocarpaceae, Cephalotaxaceae, Taxaceae, and one Cupressaceae genus (Juniperus), the scales are soft, fleshy, sweet, and brightly colored, and are eaten by fruit-eating birds, which then pass the seeds in their droppings. These fleshy scales are (except in Juniperus) known as arils. In some of these conifers (e.g. most Podocarpaceae), the cone consists of several fused scales, while in others (e.g. Taxaceae), the cone is reduced to just one seed scale or (e.g. Cephalotaxaceae) the several scales of a cone develop into individual arils, giving the appearance of a cluster of berries. The male cones have structures called microsporangia that produce yellowish pollen through meiosis. Pollen is released and carried by the wind to female cones. Pollen grains from living pinophyte species produce pollen tubes, much like those of angiosperms. The gymnosperm male gametophytes (pollen grains) are carried by wind to a female cone and are drawn into a tiny opening on the ovule called the micropyle. It is within the ovule that pollen-germination occurs. From here, a pollen tube seeks out the female gametophyte, which contains archegonia each with an egg, and if successful, fertilization occurs. The resulting zygote develops into an embryo, which along with the female gametophyte (nutritional material for the growing embryo) and its surrounding integument, becomes a seed. Eventually, the seed may fall to the ground and, if conditions permit, grow into a new plant. In forestry, the terminology of flowering plants has commonly though inaccurately been applied to cone-bearing trees as well. The male cone and unfertilized female cone are called male flower and female flower, respectively. After fertilization, the female cone is termed fruit, which undergoes ripening (maturation). It was found recently that the pollen of conifers transfers the mitochondrial organelles to the embryo, a sort of meiotic drive that perhaps explains why Pinus and other conifers are so productive, and perhaps also has bearing on observed sex-ratio bias. Life cycle Conifers are heterosporous, generating two different types of spores: male microspores and female megaspores. These spores develop on separate male and female sporophylls on separate male and female cones. In the male cones, microspores are produced from microsporocytes by meiosis. The microspores develop into pollen grains, which contain the male gametophytes. Large amounts of pollen are released and carried by the wind. Some pollen grains will land on a female cone for pollination. The generative cell in the pollen grain divides into two haploid sperm cells by mitosis leading to the development of the pollen tube. At fertilization, one of the sperm cells unites its haploid nucleus with the haploid nucleus of an egg cell. The female cone develops two ovules, each of which contains haploid megaspores. A megasporocyte is divided by meiosis in each ovule. Each winged pollen grain is a four celled male gametophyte. Three of the four cells break down leaving only a single surviving cell which will develop into a female multicellular gametophyte. The female gametophytes grow to produce two or more archegonia, each of which contains an egg. Upon fertilization, the diploid egg will give rise to the embryo, and a seed is produced. The female cone then opens, releasing the seeds which grow to a young seedling. To fertilize the ovum, the male cone releases pollen that is carried in the wind to the female cone. This is pollination. (Male and female cones usually occur on the same plant.) The pollen fertilizes the female gamete (located in the female cone). Fertilization in some species does not occur until 15 months after pollination. A fertilized female gamete (called a zygote) develops into an embryo. A seed develops which contains the embryo. The seed also contains the integument cells surrounding the embryo. This is an evolutionary characteristic of the Spermatophyta. Mature seed drops out of cone onto the ground. Seed germinates and seedling grows into a mature plant. When the plant is mature, it produces cones and the cycle continues. Female reproductive cycles Conifer reproduction is synchronous with seasonal changes in temperate zones. Reproductive development slows to a halt during each winter season and then resumes each spring. The male strobilus development is completed in a single year. Conifers are classified by three reproductive cycles that refer to the completion of female strobilus development from initiation to seed maturation. All three types of reproductive cycle have a long gap between pollination and fertilization. One year reproductive cycle: The genera include Abies, Picea, Cedrus, Pseudotsuga, Tsuga, Keteleeria (Pinaceae) and Cupressus, Thuja, Cryptomeria, Cunninghamia and Sequoia (Cupressaceae). Female strobili are initiated in late summer or fall of a year, then they overwinter. Female strobili emerge followed by pollination in the following spring. Fertilization takes place in summer of the following year, only 3–4 months after pollination. Cones mature and seeds are then shed by the end of that same year. Pollination and fertilization occur in a single growing season. Two-year reproductive cycle: The genera includes Widdringtonia, Sequoiadendron (Cupressaceae) and most species of Pinus. Female strobilus initials are formed in late summer or fall then overwinter. Female strobili emerge and receive pollen in the first year spring and become conelets. The conelet goes through another winter rest and, in the spring of the second year archegonia form in the conelet. Fertilization of the archegonia occurs by early summer of the second year, so the pollination-fertilization interval exceeds a year. After fertilization, the conelet is considered an immature cone. Maturation occurs by autumn of the second year, at which time seeds are shed. In summary, the one-year and the two-year cycles differ mainly in the duration of the pollination-fertilization interval. Three-year reproductive cycle: Three of the conifer species are pine species (Pinus pinea, Pinus leiophylla, Pinus torreyana) which have pollination and fertilization events separated by a two-year interval. Female strobili initiated during late summer or autumn of a year, then overwinter until the following spring. Female strobili emerge then pollination occurs in spring of the second year then the pollinated strobili become conelets in the same year (i.e. the second year). The female gametophytes in the conelet develop so slowly that the megaspore does not go through free-nuclear divisions until autumn of the third year. The conelet then overwinters again in the free-nuclear female gametophyte stage. Fertilization takes place by early summer of the fourth year and seeds mature in the cones by autumn of the fourth year. Tree development The growth and form of a forest tree are the result of activity in the primary and secondary meristems, influenced by the distribution of photosynthate from its needles and the hormonal gradients controlled by the apical meristems. External factors also influence growth and form. Fraser recorded the development of a single white spruce tree from 1926 to 1961. Apical growth of the stem was slow from 1926 through 1936 when the tree was competing with herbs and shrubs and probably shaded by larger trees. Lateral branches began to show reduced growth and some were no longer in evidence on the 36-year-old tree. Apical growth totaling about 340 m, 370 m, 420 m, 450 m, 500 m, 600 m, and 600 m was made by the tree in the years 1955 through 1961, respectively. The total number of needles of all ages present on the 36-year-old tree in 1961 was 5.25 million weighing 14.25 kg. In 1961, needles as old as 13 years remained on the tree. The ash weight of needles increased progressively with age from about 4% in first-year needles in 1961 to about 8% in needles 10 years old. In discussing the data obtained from the one 11 m tall white spruce, Fraser et al. (1964) speculated that if the photosynthate used in making apical growth in 1961 was manufactured the previous year, then the 4 million needles that were produced up to 1960 manufactured food for about 600,000 mm of apical growth or 730 g dry weight, over 12 million mm3 of wood for the 1961 annual ring, plus 1 million new needles, in addition to new tissue in branches, bark, and roots in 1960. Added to this would be the photosynthate to produce energy to sustain respiration over this period, an amount estimated to be about 10% of the total annual photosynthate production of a young healthy tree. On this basis, one needle produced food for about 0.19 mg dry weight of apical growth, 3 mm3 wood, one-quarter of a new needle, plus an unknown amount of branch wood, bark and roots. The order of priority of photosynthate distribution is probably: first to apical growth and new needle formation, then to buds for the next year's growth, with the cambium in the older parts of the branches receiving sustenance last. In the white spruce studied by Fraser et al. (1964), the needles constituted 17.5% of the over-day weight. Undoubtedly, the proportions change with time. Seed-dispersal mechanism Wind and animal dispersals are two major mechanisms involved in the dispersal of conifer seeds. Wind-born seed dispersal involves two processes, namely; local neighborhood dispersal and long-distance dispersal. Long-distance dispersal distances range from from the source. Birds of the crow family, Corvidae, are the primary distributor of the conifer seeds. These birds are known to cache 32,000 pine seeds and transport the seeds as far as from the source. The birds store the seeds in the soil at depths of under conditions which favor germination. Distribution and habitat Conifers are the dominant plants over large areas of land, most notably the taiga of the Northern Hemisphere, but also in similar cool climates in mountains further south. Ecology As an invasive species A number of conifers originally introduced for forestry have become invasive species in parts of New Zealand, including radiata pine (Pinus radiata), lodgepole pine (P. contorta), Douglas fir (Pseudotsuga mensiezii) and European larch (Larix decidua). In parts of South Africa, maritime pine (Pinus pinaster), patula pine (P. patula) and radiata pine have been declared invasive species. These wilding conifers are a serious environmental issue causing problems for pastoral farming and for conservation. Radiata pine was introduced to Australia in the 1870s. It is "the dominant tree species in the Australian plantation estate" – so much so that many Australians are concerned by the resulting loss of native wildlife habitat. The species is widely regarded as an environmental weed across southeastern and southwestern Australia and the removal of individual plants beyond plantations is encouraged. Predators At least 20 species of roundheaded borers of the family Cerambycidae feed on the wood of spruce, fir, and hemlock (Rose and Lindquist 1985). Borers rarely bore tunnels in living trees, although when populations are high, adult beetles feed on tender twig bark, and may damage young living trees. One of the most common and widely distributed borer species in North America is the whitespotted sawyer (Monochamus scutellatus). Adults are found in summer on newly fallen or recently felled trees chewing tiny slits in the bark in which they lay eggs. The eggs hatch in about two weeks and the tiny larvae tunnel to the wood and score its surface with their feeding channels. With the onset of cooler weather, they bore into the wood, making oval entrance holes and tunnelling deeply. Feeding continues the following summer when larvae occasionally return to the surface of the wood and extend the feeding channels generally in a U-shaped configuration. During this time, small piles of frass extruded by the larvae accumulate under logs. Early in the spring of the second year following egg-laying, the larvae, about 30 mm long, pupate in the tunnel enlargement just below the wood surface. The resulting adults chew their way out in early summer, leaving round exit holes, so completing the usual 2-year life cycle. Cultivation Conifers – notably Abies (fir), Cedrus, Chamaecyparis lawsoniana (Lawson's cypress), Cupressus (cypress), juniper, Picea (spruce), Pinus (pine), Taxus (yew), Thuja (cedar) – have been the subject of selection for ornamental purposes. Plants with unusual growth habits, sizes, and colours are propagated and planted in parks and gardens throughout the world. Conditions for growth Conifers can absorb nitrogen in either the ammonium (NH4+) or nitrate (NO3−) form, but the forms are not physiologically equivalent. Form of nitrogen affected both the total amount and relative composition of the soluble nitrogen in white spruce tissues (Durzan and Steward). Ammonium nitrogen was shown to foster arginine and amides and lead to a large increase of free guanidine compounds, whereas in leaves nourished by nitrate as the sole source of nitrogen guanidine compounds were less prominent. Durzan and Steward noted that their results, drawn from determinations made in late summer, did not rule out the occurrence of different interim responses at other times of the year. Ammonium nitrogen produced significantly heavier (dry weight) seedlings with a higher nitrogen content after 5 weeks than did the same amount of nitrate nitrogen. Swan found the same effect in 105-day-old white spruce. The general short-term effect of nitrogen fertilization on coniferous seedlings is to stimulate shoot growth more so than root growth (Armson and Carman 1961). Over a longer period, root growth is also stimulated. Many nursery managers were long reluctant to apply nitrogenous fertilizers late in the growing season, for fear of increased danger of frost damage to succulent tissues. A presentation at the North American Forest Tree Nursery Soils Workshop at Syracuse in 1980 provided strong contrary evidence: Bob Eastman, President of the Western Maine Forest Nursery Co. stated that for 15 years he has been successful in avoiding winter "burn" to Norway spruce and white spruce in his nursery operation by fertilizing with 50–80 lb/ac (56–90 kg/ha) nitrogen in September, whereas previously winter burn had been experienced annually, often severely. Eastman also stated that the overwintering storage capacity of stock thus treated was much improved (Eastman 1980). The concentrations of nutrients in plant tissues depend on many factors, including growing conditions. Interpretation of concentrations determined by analysis is easy only when a nutrient occurs in excessively low or occasionally excessively high concentration. Values are influenced by environmental factors and interactions among the 16 nutrient elements known to be essential to plants, 13 of which are obtained from the soil, including nitrogen, phosphorus, potassium, calcium, magnesium, and sulfur, all used in relatively large amounts. Nutrient concentrations in conifers also vary with season, age, and kind of tissue sampled, and analytical technique. The ranges of concentrations occurring in well-grown plants provide a useful guide by which to assess the adequacy of particular nutrients, and the ratios among the major nutrients are helpful guides to nutritional imbalances. Economic importance The softwood derived from conifers is of great economic value, providing about 45% of the world's annual lumber production. Other uses of the timber include the production of paper and plastic from chemically treated wood pulp. Some conifers also provide foods such as pine nuts and juniper berries, the latter used to flavor gin.	Biology and health sciences	Gymnosperms	null
68094	https://en.wikipedia.org/wiki/Pinales	Pinales	The order Pinales in the division Pinophyta, class Pinopsida, comprises all the extant conifers. The distinguishing characteristic is the reproductive structure known as a cone produced by all Pinales. All of the extant conifers, such as Araucaria, cedar, celery-pine, cypress, fir, juniper, kauri, larch, pine, redwood, spruce, and yew, are included here. Some fossil conifers, however, belong to other distinct orders within the division Pinophyta. Multiple molecular studies indicate this order being paraphyletic with respect to Gnetales, with studies recovering Gnetales as either a sister group to Pinaceae or being more derived than Pinaceae but sister to the rest of the group. Taxonomy History Brown (1825) first discerned that there were two groups of seed plants, distinguished by the form of seed development, based on whether the ovules were exposed, receiving pollen directly, or enclosed, which do not. Shortly afterwards, Brongniart (1828) coined the term Phanérogames gymnosperms to describe the former group. The distinction was then formalized by Lindley (1830), dividing what he referred to as the subclass Dicotyledons into two tribes, Gymnosperms and Angiosperms. In the gymnosperms (or Gymnospermae) Lindley included two orders, the Cycadeae and the Coniferae. In his final work (1853) he described Gymnogens as a class with four orders; Cycadeaceae (cycads) Pinaceae (conifers) Taxaceae (taxads) Gnetaceae In contrast, Bentham and Hooker (1880) included only three orders in the class Gymnospermeae, by including taxads within Coniferae; Gnetaceae Coniferae Cycadaceae In the Engler system (1903) Gymnospermae is listed as a subdivision (Unterabteilung) and adopted more of a splitter approach, including extinct taxa, with the following six classes; Cycadales Bennettitales Cordaitales Ginkgoales Coniferae Gnetales During this period, Gorozhankin published his treatise on Gymnosperms (1895), for which he bears the botanical authority for Pinales, Gorozh.. In his classification, Gymnospermae (alternatively named Archespermae) was a class of the division Archegoniatae, divided into subclasses; Cycadoideae Peucideae (Coniferae) A system of two groups was maintained by the most commonly used classification in the twentieth century, the revision of the Engler system by Pilger (1926), who grouped 12 families of the Gymnospermae subdivision into 2 classes; Coniferales (Coniferae) Gnetales The treatment of Gymnosperms as two groups, though with varying composition and names, was followed for most of the twentieth century, including the systems of Chamberlain (1935), Benson (1957) and Cronquist (1960). In the latter, Cronquist divided Gymnospermae into two divisions; division Coniferophyta class Coniferae class Chlamydospermae (Gnetales) division Cycadophyta class Cycadae Benson,(1957) who introduced the term Pinales, divided gymnosperms into four classes; Conopsida (conifers, including Pinales) Ephedropsida Gnetopsida Cycadopsida In a later revision, in collaboration with two other taxonomists (1966), Cronquist merged all the gymnosperms into a single division, Pinophyta, with three subdivisions reflecting the main lineages; Cycadicae Pinicae Gneticae In the era of molecular phylogenetics, De-Zhi and colleagues (2004) once again proposed a division of 12 gymnosperm families into two classes; Cycadopsida order Cycadales Coniferopsida subclass Multinervidae (6 orders) subclass Taxidae order Taxales order Pinales With the development of the Angiosperm Phylogeny Group came a major realignment of the linear classification of the land plants, by Chase and Reveal (2009). In this system, the land plants form a class, Equisetopsida s.l. (sensu lato) or sensu Chase & Reveal, also known as embryophytes or Embryophyceae nom. illeg.. Class Equisetopsida s.l. is divided into 14 subclades as subclasses, including Magnoliidae (angiosperms). The gymnosperms are represented by four of these subclasses, placing them in a sister group relationship to angiosperms. Subclasses (number of orders); subclass Cycadidae Pax (1) subclass Ginkgooidae Engl. (1) subclass Gnetidae Pax (3) subclass Pinidae Cronquist, Takht. & Zimmerm. (conifers) (1) Controversies Gymnosperm (Acrogymnospermae) taxonomy has been considered controversial, and lacks consensus. As taxonomic classification transformed from being based solely on plant morphology to molecular phylogenetics, the number of taxonomic publications increased considerably after 2008, however, these approaches have not been uniform. A taxonomic classification has been complicated by the relationship of extant to extinct taxa, and within extinct taxa, and particularly the placement of Gnetophyta. The latter have been variously classified as basal to all gymnosperms, sister group to conifers (‘gnetifer’ hypothesis) or sister to Pinaceae (‘gnepine’ hypothesis) in which they are classified within the conifers. The extant conifers most likely form a monophyletic group. In 2018, the Gymnosperm Phylogeny Group was established, analogous to the Angiosperm Phylogeny Group and Pteridophyte Phylogeny Group, with the intention of reaching a consensus. Phylogeny Gymnosperms form a group of four subclasses among the spermatophytes (seed bearing plants). In turn, the seed plants together with the monilophyte fern subclasses make up the tracheophytes (vascular plants), part of the class Equisetopsida (embryophytes or land plants), as opposed to the green algae. Among the seed plants, the gymnosperms are a sister group to the subclass Magnoliidae (angiosperms or flowering plants). There are about 1000 extant gymnosperm species, distributed over about 12 families and 83 genera. Many of these genera are monotypic (41%), and another 27% are oligotypic (2–5 species). The four subclasses have also been treated as divisions of the Spermatophytes. Alternative names and the approximate number of genera and species in each are; Cycadidae (Cycadophyta, cycads 10, 300) Ginkgoidae (Ginkgophyta, ginkgo 1, 1) Gnetidae (Gnetophyta 3, 100) Pinidae (Pinophyta, conifers 70, 600) The term Pinophyta has also been used to include all conifers, extinct and extant, with Pinales representing all the extant conifers. Christenhusz and colleagues extended the system of Chase and Reveal to provide a revised classification of gymnosperms in 2011, based on the above four subclades. In this scheme, the Pinidae comprise three orders, including Pinales, and 6 families; Pinales Gorozh. (Pinaceae) Araucariales Gorozh. (Araucariaceae, Podocarpaceae) Cupressales Link (Sciadopityaceae, Cupressaceae, Taxaceae) However, the exact phylogeny remained a topic that was 'hotly debated", in particular whether the main lineages were best represented by the four subclasses of Christenhusz and colleagues or the more traditional five clades (cycads, ginkgos, cupressophytes, Pinaceae and gnetophytes). In 2014 the first complete molecular phylogeny was published, based on 90 species representing all extant genera. This established cycads as the basal group, followed by Ginkgoaceae, as sister to the remaining gymnosperms, and supporting the ‘gnepine’ hypothesis. This analysis favours the five clade hypothesis, the remaining clade following divergence of the Pinidae, are referred to as the conifer II clade, or cupressophytes, in distinction from the conifer I clade (Gnetidae, Pinidae). This clade, in turn, has two lineages. The first consisting of Sciadopityaceae and the Araucariales, the second being the Cupressales. In the Christenhusz scheme, the Sciadopityaceae were considered to be within Cupressales. The term Cupressaceae s.l. refers to the inclusion of Taxodiaceae. These relationships are shown in this cladogram, although no formal taxonomic revision was undertaken. A more comprehensive analysis was undertaken by Ran and colleagues in 2018, as part of a detailed phylogeny of all seed plants. This forms the basis of the Tracheophyte Phylogeny Poster and the Angiosperm Phylogeny Website. Subdivision Historically conifers, in the order Pinales have been considered to consist of six to seven extant families, based on the classification of class Coniferae by Pilger (1926), considered the standard through most of the twentieth century. These families were treated as a single order, in distinction to some earlier systems. His families were; Araucariaceae Cupressaceae (cypresses, juniper, redwood) Pinaceae (firs, pines, cedars, larch, spruce) Podocarpaceae Taxaceae (yews) Cephalotaxaceae Taxodiaceae Subsequent revisions merged the Taxodiaceae and Cupressaceae, and placed Sciadopitys, formerly in Cupressaceae, into a separate family (Sciadopityaceae). Cephalotaxaceae had previously been recognized as a separate family, but was subsequently included in Taxaceae. Similarly Phyllocladaceae were included in Podocarpaceae. Yews (Taxaceae) have sometimes been treated as a separate order (Taxales). Christenhusz and colleagues (2011) included only one family in Pinales, Pinaceae, a practice subsequently followed by the Angiosperm Phylogeny Website and the Gymnosperm Database. In this restricted model Pinales (Pinaceae) comprisea 11 genera and about 225 species, all of the other conifers originally included in this order, being included in other orders such as Cupressales.	Biology and health sciences	Pinophyta (Conifers)	Plants
68110	https://en.wikipedia.org/wiki/Logging	Logging	Logging is the process of cutting, processing, and moving trees to a location for transport. It may include skidding, on-site processing, and loading of trees or logs onto trucks or skeleton cars. In forestry, the term logging is sometimes used narrowly to describe the logistics of moving wood from the stump to somewhere outside the forest, usually a sawmill or a lumber yard. In common usage, however, the term may cover a range of forestry or silviculture activities. Logging is the beginning of a supply chain that provides raw material for many products societies worldwide use for housing, construction, energy, and consumer paper products. Logging systems are also used to manage forests, reduce the risk of wildfires, and restore ecosystem functions, though their efficiency for these purposes has been challenged. Logging frequently has negative impacts. The harvesting procedure itself may be illegal, including the use of corrupt means to gain access to forests; extraction without permission or from a protected area; the cutting of protected species; or the extraction of timber in excess of agreed limits. It may involve the so-called "timber mafia". Excess logging can lead to irreparable harm to ecosystems, such as deforestation and biodiversity loss. Infrastructure for logging can also lead to other environmental degradation. These negative environmental impacts can lead to environmental conflict. Additionally, there is significant occupational injury risk involved in logging. Logging can take many formats. Clearcutting (or "block cutting") is not necessarily considered a type of logging but a harvesting or silviculture method. Cutting trees with the highest value and leaving those with lower value, often diseased or malformed trees, is referred to as high grading. It is sometimes called selective logging, and confused with selection cutting, the practice of managing stands by harvesting a proportion of trees. Logging usually refers to above-ground forestry logging. Submerged forests exist on land that has been flooded by damming to create reservoirs. Harvesting trees from forests submerged by flooding or dam creation is called underwater logging, a form of timber recovery. Clearcutting Clearcutting, or clearfelling, is a method of harvesting that removes essentially all the standing trees in a selected area. Depending on management objectives, a clearcut may or may not have reserve trees left to attain goals other than regeneration, including wildlife habitat management, mitigation of potential erosion or water quality concerns. Silviculture objectives for clearcutting, (for example, healthy regeneration of new trees on the site) and a focus on forestry distinguish it from deforestation. Other methods include shelterwood cutting, group selective, single selective, seed-tree cutting, patch cut, and retention cutting. Logging methods The above operations can be carried out by different methods, of which the following three are considered industrial methods: / stem-only harvesting Trees are felled and then delimbed and topped at the stump. The log is then transported to the landing, where it is bucked and loaded on a truck. This leaves the slash (and the nutrients it contains) in the cut area, where it must be further treated if wild land fires are of concern. Whole-tree logging Trees and plants are felled and transported to the roadside with top and limbs intact. There have been advancements to the process which now allows a logger or harvester to cut the tree down, top, and delimb a tree in the same process. This ability is due to the advancement in the style felling head that can be used. The trees are then delimbed, topped, and bucked at the landing. This method requires that slash be treated at the landing. In areas with access to cogeneration facilities, the slash can be chipped and used for the production of electricity or heat. Full-tree harvesting also refers to utilization of the entire tree including branches and tops. This technique removes both nutrients and soil cover from the site and so can be harmful to the long-term health of the area if no further action is taken, however, depending on the species, many of the limbs are often broken off in handling so the result may not be as different from tree-length logging as it might seem. Cut-to-length logging Cut-to-length logging is the process of felling, delimbing, bucking, and sorting (pulpwood, sawlog, etc.) at the stump area, leaving limbs and tops in the forest. Mechanical harvesters fell the tree, delimb, and buck it, and place the resulting logs in bunks to be brought to the landing by a skidder or forwarder. This method is routinely available for trees up to in diameter. Transporting logs Logging methods have changed over time, driven by advancements in transporting timber from remote areas to markets. These shifts fall into three main eras: the manual logging era before the 1880s, the railroad logging era from the 1880s to World War II, and the modern mechanized era that began after the war. Pre-1880s: Pre-Industrial Era In the early days, felled logs were transported using simple methods such as rivers to float tree trunks downstream to sawmills or paper mills. This practice, known as log driving or timber rafting, was the cheapest and most common. Some logs, due to high resin content, would sink and were known as deadheads. Logs were also moved with high-wheel loaders, a set of wheels over ten feet tall, initially pulled by oxen. 1880s to World War II: Railroad Logging Era As the logging industry expanded, the 1880s saw the introduction of mechanized equipment like railroads and steam-powered machinery, marking the beginning of the railroad logging era. Logs were moved more efficiently by railroads built into remote forest areas, often supported by additional methods like high-wheel loaders, tractors and log flumes. The largest high-wheel loader, the "Bunyan Buggie," was built in 1960 for service in California, featuring wheels high. Post-World War II: Modern Mechanized Logging After World War II, mechanized logging equipment, including chainsaws, diesel trucks, and Caterpillar tractors, transformed the logging industry, making railroad-based logging obsolete. With the advent of these tools, transporting logs became more efficient as new roads were constructed to access remote forests. However, in protected areas like United States National Forests and designated wilderness zones, road building has been restricted to minimize environmental impacts such as erosion in riparian zones. Today, heavy machinery such as yarders and skyline systems are used to gather logs from steep terrain, while helicopters are used for heli-logging to minimize environmental impact. Less common forms of logging, like horse logging and the use of oxen, still exist but are mostly superseded. Safety considerations Logging is a dangerous occupation. In the United States, it has consistently been one of the most hazardous industries and was recognized by the National Institute for Occupational Safety and Health (NIOSH) as a priority industry sector in the National Occupational Research Agenda (NORA) to identify and provide intervention strategies regarding occupational health and safety issues. In 2008, the logging industry employed 86,000 workers and accounted for 93 deaths. This resulted in a fatality rate of 108.1 deaths per 100,000 workers that year. This rate is over 30 times higher than the overall fatality rate. Forestry/logging-related injuries (fatal and non-fatal) are often difficult to track through formal reporting mechanisms. Thus, some programs have begun to monitor injuries through publicly available reports such as news media. The logging industry experiences the highest fatality rate of 23.2 per 100,000 full-time equivalent (FTE) workers and a non-fatal incident rate of 8.5 per 100 FTE workers. The most common type of injuries or illnesses at work include musculoskeletal disorders (MSDs), which include an extensive list of "inflammatory and degenerative conditions affecting the muscles, tendons, ligaments, joints, peripheral nerves, and supporting blood vessels." Loggers work with heavy, moving weights, and use tools such as chainsaws and heavy equipment on uneven and sometimes steep or unstable terrain. Loggers also deal with severe environmental conditions, such as inclement weather and severe heat or cold. An injured logger is often far from professional emergency treatment. Traditionally, the cry of "Timber!" developed as a warning alerting fellow workers in an area that a tree is being felled, so they should be alert to avoid being struck. The term "widowmaker" for timber, typically a limb or branch that is no longer attached to a tree, but is still in the canopy either wedged in a crotch, tangled in other limbs, or miraculously balanced on another limb demonstrates another emphasis on situational awareness as a safety principle. In British Columbia, Canada, the BC Forest Safety Council was created in September 2004 as a not-for-profit society dedicated to promoting safety in the forest sector. It works with employers, workers, contractors, and government agencies to implement fundamental changes necessary to make it safer to earn a living in forestry. The risks experienced in logging operations can be somewhat reduced, where conditions permit, by the use of mechanical tree harvesters, skidders, and forwarders.	Technology	Trees and forestry	null
68181	https://en.wikipedia.org/wiki/Workstation	Workstation	A workstation is a special computer designed for technical or scientific applications. Intended primarily to be used by a single user, they are commonly connected to a local area network and run multi-user operating systems. The term workstation has been used loosely to refer to everything from a mainframe computer terminal to a PC connected to a network, but the most common form refers to the class of hardware offered by several current and defunct companies such as Sun Microsystems, Silicon Graphics, Apollo Computer, DEC, HP, NeXT, and IBM which powered the 3D computer graphics revolution of the late 1990s. Workstations formerly offered higher performance than mainstream personal computers, especially in CPU, graphics, memory, and multitasking. Workstations are optimized for the visualization and manipulation of different types of complex data such as 3D mechanical design, engineering simulations like computational fluid dynamics, animation, video editing, image editing, medical imaging, image rendering, computational science, generating mathematical plots, and software development. Typically, the form factor is that of a desktop computer, which consists of a high-resolution display, a keyboard, and a mouse at a minimum, but also offers multiple displays, graphics tablets, and 3D mice for manipulating objects and navigating scenes. Workstations were the first segment of the computer market to present advanced accessories, and collaboration tools like videoconferencing. The increasing capabilities of mainstream PCs since the late 1990s have reduced distinction between the PCs and workstations. Typical 1980s workstations have expensive proprietary hardware and operating systems to categorically distinguish from standardized PCs. From the 1990s and 2000s, IBM's RS/6000 and IntelliStation have RISC-based POWER CPUs running AIX, versus its corporate IBM PC Series and consumer Aptiva PCs that have Intel x86 CPUs and usually running Microsoft Windows. However, by the early 2000s, this difference largely disappeared, since workstations use highly commoditized hardware dominated by large PC vendors, such as Dell, Hewlett-Packard, and Fujitsu, selling x86-64 systems running Windows or Linux. History Origins and development Workstations are older than the first personal computer (PC). The first computer that might qualify as a workstation is the IBM 1620, a small scientific computer designed to be used interactively by a single person sitting at the console. It was introduced in 1959. One peculiar feature of the machine is that it lacks any arithmetic circuitry. To perform addition, it requires a memory-resident table of decimal addition rules. This reduced the cost of logic circuitry, enabling IBM to make it inexpensive. The machine is codenamed CADET and was initially rented for $1000 per month. In 1965, the IBM 1130 scientific computer became the successor to 1620. Both of these systems run Fortran and other languages. They are built into roughly desk-sized cabinets, with console typewriters. They have optional add-on disk drives, printers, and both paper-tape and punched-card I/O. Early workstations were generally dedicated minicomputers, a multiuser system reserved for one user. For example, the PDP-8 from Digital Equipment Corporation, is regarded as the first commercial minicomputer. Workstations have historically been more advanced than contemporary PCs, with more powerful CPU architectures, earlier networking, more advanced graphics, more memory, and multitasking with sophisticated operating systems like Unix. Because of their minicomputer heritage, from the start workstations have run professional and expensive software such as CAD and graphics design, as opposed to PCs' games and text editors. The Lisp machines developed at MIT in the early 1970s pioneered some workstation principles, as high-performance, networked, single-user systems intended for heavily interactive use. Lisp Machines were commercialized beginning 1980 by companies like Symbolics, Lisp Machines, Texas Instruments (the TI Explorer), and Xerox (the Interlisp-D workstations). The first computer designed for a single user, with high-resolution graphics (and so a workstation in the modern sense), is the Alto developed at Xerox PARC in 1973. Other early workstations include the Terak 8510/a (1977), Three Rivers PERQ (1979), and the later Xerox Star (1981). 1980s rise in popularity In the early 1980s, with the advent of 32-bit microprocessors such as the Motorola 68000, several new competitors appeared, including Apollo Computer and Sun Microsystems, with workstations based on 68000 and Unix. Meanwhile, DARPA's VLSI Project created several spinoff graphics products, such as the Silicon Graphics 3130. Target markets were differentiated, with Sun and Apollo considered to be network workstations and SGI as graphics workstations. RISC CPUs increased in the mid-1980s, typical of workstation vendors. Competition between RISC vendors lowered CPU prices to as little as $10 per MIPS, much less expensive than the Intel 80386; after large price cuts in 1987 and 1988, a personal workstation suitable for 2D CAD costing to was available from multiple vendors. Mid-range models capable of 3D graphics cost from to , while high-end models overlapping with minicomputers cost from to or more. By then a "personal workstation" might be a high-end PC like Macintosh II or IBM PS/2 Model 80, low-end workstation, or a hybrid device like the NeXT Computer, all with similar, overlapping specifications. One differentiator between PC and workstation was that the latter was much more likely to have a graphics accelerator with support for a graphics standard like PHIGS or X Window, while the former usually depended on software rendering or proprietary accelerators. The computer animation industry's needs typically caused improvements in graphical technology, with CAD using the same improvements later. BYTE predicted in 1989 "Soon, the only way we'll be able to tell the difference between traditional workstations and PCs will be by the operating system they run", with the former running Unix and the latter running OS/2, classic Mac OS, and/or Unix. Many workstations by then had some method to run increasingly popular and powerful PC software such as Lotus 1-2-3 or Microsoft Word. The magazine demonstrated that year that an individual could build a workstation with commodity components with specifications comparable to commercially available low-end workstations. Workstations often featured SCSI or Fibre Channel disk storage systems, high-end 3D accelerators, single or multiple 64-bit processors, large amounts of RAM, and well-designed cooling. Additionally, the companies that make the products tend to have comprehensive repair/replacement plans. As the distinction between workstation and PC fades, however, workstation manufacturers have increasingly employed "off-the-shelf" PC components and graphics solutions rather than proprietary hardware or software. Some "low-cost" workstations are still expensive by PC standards but offer binary compatibility with higher-end workstations and servers made by the same vendor. This allows software development to take place on low-cost (relative to the server) desktop machines. Thin clients Workstations diversified to the lowest possible price point as opposed to performance, called the thin client or network computer. Dependent upon a network and server, this reduces the machine to having no hard drive, and only the CPU, keyboard, mouse, and screen. Some diskless nodes still run a traditional operating system and perform computations locally, with storage on a remote server. These are intended to reduce the initial system purchase cost, and the total cost of ownership, by reducing the amount of administration required per user. This approach was first attempted as a replacement for PCs in office productivity applications, with the 3Station by 3Com. In the 1990s, X terminals filled a similar role for technical computing. Sun's thin clients include the Sun Ray product line. However, traditional workstations and PCs continued to drop in price and complexity as remote management tools for IT staff became available, undercutting this market. 3M computer A high-end workstation of the early 1980s with the three Ms, or a "3M computer" (coined by Raj Reddy and his colleagues at CMU), has one megabyte of RAM, a megapixel display (roughly 1000×1000 pixels), and one "MegaFLOPS" compute performance (at least one million floating-point operations per second). RFC 782 defines the workstation environment more generally as "hardware and software dedicated to serve a single user", and that it provisions additional shared resources. This is at least one order of magnitude beyond the capacity of the personal computer of the time. The original 1981 IBM Personal Computer has 16 KB memory, a text-only display, and floating-point performance around ( with the optional 8087 math coprocessor. Other features beyond the typical personal computer include networking, graphics acceleration, and high-speed internal and peripheral data buses. Another goal was to bring the price below one "megapenny", that is, less than , which was achieved in the late 1980s. Throughout the early to mid-1990s, many workstations cost from to or more. Decline The more widespread adoption of these technologies into mainstream PCs was a direct factor in the decline of the workstation as a separate market segment: Reliable components High-performance 3D graphics hardware for computer-aided design (CAD) and computer-generated imagery (CGI) animation is increasingly popular in the PC market around the mid-to-late 1990s mostly driven by computer gaming, yielding the first official GPU in Nvidia's NV10 and the breakthrough GeForce 256. High-performance CPUs: the first RISC of the early 1980s offer roughly one order of magnitude in performance improvement over CISC processors of comparable cost. Intel's x86 CISC family always had the edge in market share and the economies of scale that this implied. By the mid-1990s, some CISC processors like the Motorola 68040 and Intel's 80486 and Pentium have performance parity with RISC in some areas, such as integer performance (at the cost of greater chip complexity) and hardware floating-point calculations, relegating RISC to even more high-end markets. Hardware support for floating-point operations: optional on the original IBM PC; remained on a separate chip for Intel systems until the 80486DX processor. Even then, x86 floating-point performance lags other processors due to limitations in its architecture. Today even low-price PCs now have performance in the gigaFLOPS range. High-performance/high-capacity data storage: early workstations tend to use proprietary disk interfaces until the SCSI standard of the mid-1980s. Although SCSI interfaces soon became available for IBM PCs, they were comparatively expensive and tend to be limited by the speed of the PC's ISA peripheral bus. SCSI is an advanced controller interface good for multitasking and daisy chaining. This makes it suited for use in servers, and its benefits to desktop PCs which mostly run single-user operating systems are less clear, but it is standard on the 1980s-1990s Macintosh. Serial ATA is more modern, with throughput comparable to SCSI but at a lower cost. High-speed networking (10 Mbit/s or better): 10 Mbit/s network interfaces were commonly available for PCs by the early 1990s, although by that time workstations were pursuing even higher networking speeds, moving to 100 Mbit/s, 1 Gbit/s, and 10 Gbit/s. However, economies of scale and the demand for high-speed networking in even non-technical areas have dramatically decreased the time it takes for newer networking technologies to reach commodity price points. Large displays (17- to 21-inch) with high resolutions and high refresh rates for graphics and CAD work, which were rare among PCs in the late 1980s and early 1990s but became common among PCs by the late 1990s. Large memory configurations: PCs (such as IBM clones) are originally limited to 640 KB of RAM until the 1982 introduction of the 80286 processor; early workstations have megabytes of memory. IBM clones require special programming techniques to address more than 640 KB until the 80386, as opposed to other 32-bit processors such as SPARC which provide straightforward access to nearly their entire 4 GB memory address range. 64-bit workstations and servers supporting an address range far beyond 4 GB have been available since the early 1990s, a technology just beginning to appear in the PC desktop and server market in the mid-2000s. Operating system: early workstations ran the Unix operating system (OS), a Unix-like variant, or an unrelated equivalent OS such as VMS. The PC CPUs of the time have limitations in memory capacity and memory access protection, making them unsuitable to run OSes of this sophistication, but this, too, began to change in the late 1980s as PCs with the 32-bit 80386 with integrated paged MMUs became widely affordable and enabling OS/2, Windows NT 3.1, and Unix-like systems based on BSD and Linux on commodity PC hardware. Tight integration between the OS and the hardware: Workstation vendors both design the hardware and maintain the Unix operating system variant that runs on it. This allows for much more rigorous testing than is possible with an operating system such as Windows. Windows requires that third-party hardware vendors write compliant hardware drivers that are stable and reliable. Also, minor variations in hardware quality such as timing or build quality can affect the reliability of the overall machine. Workstation vendors are able to ensure both the quality of the hardware, and the stability of the operating system drivers by validating these things in-house, and this leads to a generally much more reliable and stable machine. Market position Since the late 1990s, the workstation and consumer markets have further merged. Many low-end workstation components are now the same as the consumer market, and the price differential narrowed. For example, most Macintosh Quadra computers were originally intended for scientific or design work, all with the Motorola 68040 CPU, backward compatible with 68000 Macintoshes. The consumer Macintosh IIcx and Macintosh IIci models can be upgraded to the Quadra 700. "In an era when many professionals preferred Silicon Graphics workstations, the Quadra 700 was an intriguing option at a fraction of the cost" as resource-intensive software such as Infini-D brought "studio-quality 3D rendering and animations to the home desktop". The Quadra 700 can run A/UX 3.0, making it a Unix workstation. Another example is the Nvidia GeForce 256 consumer graphics card, which spawned the Quadro workstation card, which has the same GPU but different driver support and certifications for CAD applications and a much higher price. Workstations have typically driven advancements in CPU technology. All computers benefit from multi-processor and multicore designs (essentially, multiple processors on a die). The multicore design was pioneered by IBM's POWER4; it and Intel Xeon have multiple CPUs, more on-die cache, and ECC memory. Some workstations are designed or certified for use with only one specific application such as AutoCAD, Avid Xpress Studio HD, or 3D Studio Max. The certification process increases workstation prices. Modern market GPU workstations Modern workstations are typically desktop computers with AMD or NVIDIA GPUs to do high-performance computing on software programs such as video editing, 3D modeling, computer-aided design, and rendering. Decline of RISC workstations By January 2009, all RISC-based workstation product lines had been discontinued: Hewlett-Packard withdrew its last HP 9000 PA-RISC-based desktop products from the market in January 2008. IBM retired the IntelliStation POWER on January 2, 2009. SGI ended general availability of its MIPS-based SGI Fuel and SGI Tezro workstations in December 2006. Sun Microsystems announced end-of-life for its last Sun Ultra SPARC workstations in October 2008. In early 2018, RISC workstations were reintroduced in a series of IBM POWER9-based systems by Raptor Computing Systems. In October of 2024 System 76 introduces The Thelio Astra an Arm workstation aim for autonomous car industry. x86-64 Most of the current workstation market uses x86-64 microprocessors. Operating systems include Windows, FreeBSD, Linux distributions, macOS, and Solaris. Some vendors also market commodity mono-socket systems as workstations. These are three types of workstations: Workstation blade systems (IBM HC10 or Hewlett-Packard xw460c. Sun Visualization System is akin to these solutions) Ultra high-end workstation (SGI Virtu VS3xx) Deskside systems containing server-class CPUs and chipsets on large server-class motherboards with high-end RAM (HP Z-series workstations and Fujitsu CELSIUS workstations) Definition A high-end desktop market segment includes workstations, with PC operating systems and components. Component product lines may be segmented, with premium components that are functionally similar to the consumer models but with higher robustness or performance. A workstation-class PC may have some of the following features: Larger number of memory sockets which use DIMM slots or registered (buffered) modules Multiple displays Reliable high-performance graphics card Multiple processor sockets, powerful CPUs Run reliable operating system with advanced features Support for ECC memory M.2 or PCI-E NVMe SSD	Technology	Computer hardware	null
68206	https://en.wikipedia.org/wiki/Central%20dogma%20of%20molecular%20biology	Central dogma of molecular biology	The central dogma of molecular biology deals with the flow of genetic information within a biological system. It is often stated as "DNA makes RNA, and RNA makes protein", although this is not its original meaning. It was first stated by Francis Crick in 1957, then published in 1958: He re-stated it in a Nature paper published in 1970: "The central dogma of molecular biology deals with the detailed residue-by-residue transfer of sequential information. It states that such information cannot be transferred back from protein to either protein or nucleic acid." A second version of the central dogma is popular but incorrect. This is the simplistic DNA → RNA → protein pathway published by James Watson in the first edition of The Molecular Biology of the Gene (1965). Watson's version differs from Crick's because Watson describes a two-step (DNA → RNA and RNA → protein) process as the central dogma. While the dogma as originally stated by Crick remains valid today, Watson's version does not. Biological sequence information The biopolymers that comprise DNA, RNA and (poly)peptides are linear heteropolymers (i.e.: each monomer is connected to at most two other monomers). The sequence of their monomers effectively encodes information. The transfers of information from one molecule to another are faithful, deterministic transfers, wherein one biopolymer's sequence is used as a template for the construction of another biopolymer with a sequence that is entirely dependent on the original biopolymer's sequence. When DNA is transcribed to RNA, its complement is paired to it. DNA codes are transferred to RNA codes in a complementary fashion. The encoding of proteins is done in groups of three, known as codons. The standard codon table applies for humans and mammals, but some other lifeforms (including human mitochondria) use different translations. General transfers of biological sequential information DNA replications In the sense that DNA replication must occur if genetic material is to be provided for the progeny of any cell, whether somatic or reproductive, the copying from DNA to DNA arguably is the fundamental step in information transfer. A complex group of proteins called the replisome performs the replication of the information from the parent strand to the complementary daughter strand. Transcription Transcription is the process by which the information contained in a section of DNA is replicated in the form of a newly assembled piece of messenger RNA (mRNA). Enzymes facilitating the process include RNA polymerase and transcription factors. In eukaryotic cells the primary transcript is pre-mRNA. Pre-mRNA must be processed for translation to proceed. Processing includes the addition of a 5' cap and a poly-A tail to the pre-mRNA chain, followed by splicing. Alternative splicing occurs when appropriate, increasing the diversity of the proteins that any single mRNA can produce. The product of the entire transcription process (that began with the production of the pre-mRNA chain) is a mature mRNA chain. Translation The mature mRNA finds its way to a ribosome, where it gets translated. In prokaryotic cells, which have no nuclear compartment, the processes of transcription and translation may be linked together without clear separation. In eukaryotic cells, the site of transcription (the cell nucleus) is usually separated from the site of translation (the cytoplasm), so the mRNA must be transported out of the nucleus into the cytoplasm, where it can be bound by ribosomes. The ribosome reads the mRNA triplet codons, usually beginning with an AUG (adenine−uracil−guanine), or initiator methionine codon downstream of the ribosome binding site. Complexes of initiation factors and elongation factors bring aminoacylated transfer RNAs (tRNAs) into the ribosome-mRNA complex, matching the codon in the mRNA to the anti-codon on the tRNA. Each tRNA bears the appropriate amino acid residue to add to the polypeptide chain being synthesised. As the amino acids get linked into the growing peptide chain, the chain begins folding into the correct conformation. Translation ends with a stop codon which may be a UAA, UGA, or UAG triplet. The mRNA does not contain all the information for specifying the nature of the mature protein. The nascent polypeptide chain released from the ribosome commonly requires additional processing before the final product emerges. For one thing, the correct folding process is complex and vitally important. For most proteins it requires other chaperone proteins to control the form of the product. Some proteins then excise internal segments from their own peptide chains, splicing the free ends that border the gap; in such processes the inside "discarded" sections are called inteins. Other proteins must be split into multiple sections without splicing. Some polypeptide chains need to be cross-linked, and others must be attached to cofactors such as haem (heme) before they become functional. Additional transfers of biological sequential information Reverse transcription Reverse transcription is the transfer of information from RNA to DNA (the reverse of normal transcription). This is known to occur in the case of retroviruses, such as HIV, as well as in eukaryotes, in the case of retrotransposons and telomere synthesis. It is the process by which genetic information from RNA gets transcribed into new DNA. The family of enzymes involved in this process is called Reverse Transcriptase. RNA replication RNA replication is the copying of one RNA to another. Many viruses replicate this way. The enzymes that copy RNA to new RNA, called RNA-dependent RNA polymerases, are also found in many eukaryotes where they are involved in RNA silencing. RNA editing, in which an RNA sequence is altered by a complex of proteins and a "guide RNA", could also be seen as an RNA-to-RNA transfer. Activities unrelated to the central dogma The central dogma of molecular biology states that once sequential information has passed from nucleic acid to protein it cannot flow back from protein to nucleic acid. Some people believe that the following activities conflict with the central dogma. Post-translational modification After protein amino acid sequences have been translated from nucleic acid chains, they can be edited by appropriate enzymes. This is a form of protein affecting protein sequence not protein transferring information to nucleic acid. Nonribosomal peptide synthesis Some proteins are synthesized by nonribosomal peptide synthetases, which can be big protein complexes, each specializing in synthesizing only one type of peptide. Nonribosomal peptides often have cyclic and/or branched structures and can contain non-proteinogenic amino acids - both of these factors differentiate them from ribosome synthesized proteins. An example of nonribosomal peptides are some of the antibiotics. Inteins An intein is a "parasitic" segment of a protein that is able to excise itself from the chain of amino acids as they emerge from the ribosome and rejoin the remaining portions with a peptide bond in such a manner that the main protein "backbone" does not fall apart. This is a case of a protein changing its own primary sequence from the sequence originally encoded by the DNA of a gene. Additionally, most inteins contain a homing endonuclease or HEG domain which is capable of finding a copy of the parent gene that does not include the intein nucleotide sequence. On contact with the intein-free copy, the HEG domain initiates the DNA double-stranded break repair mechanism. This process causes the intein sequence to be copied from the original source gene to the intein-free gene. This is an example of protein directly editing DNA sequence, as well as increasing the sequence's heritable propagation. Prions Prions are proteins of particular amino acid sequences in particular conformations. They propagate themselves in host cells by making conformational changes in other molecules of protein with the same amino acid sequence, but with a different conformation that is functionally important or detrimental to the organism. Once the protein has been transconformed to the prion folding it changes function. In turn it can convey information into new cells and reconfigure more functional molecules of that sequence into the alternate prion form. In some types of prion in fungi this change is continuous and direct; the information flow is Protein → Protein. Some scientists such as Alain E. Bussard and Eugene Koonin have argued that prion-mediated inheritance violates the central dogma of molecular biology. However, Rosalind Ridley in Molecular Pathology of the Prions (2001) has written that "The prion hypothesis is not heretical to the central dogma of molecular biology—that the information necessary to manufacture proteins is encoded in the nucleotide sequence of nucleic acid—because it does not claim that proteins replicate. Rather, it claims that there is a source of information within protein molecules that contributes to their biological function, and that this information can be passed on to other molecules." Use of the term dogma In his autobiography, What Mad Pursuit, Crick wrote about his choice of the word dogma and some of the problems it caused him: "I called this idea the central dogma, for two reasons, I suspect. I had already used the obvious word hypothesis in the sequence hypothesis, and in addition I wanted to suggest that this new assumption was more central and more powerful. ... As it turned out, the use of the word dogma caused almost more trouble than it was worth. Many years later Jacques Monod pointed out to me that I did not appear to understand the correct use of the word dogma, which is a belief that cannot be doubted. I did apprehend this in a vague sort of way but since I thought that all religious beliefs were without foundation, I used the word the way I myself thought about it, not as most of the world does, and simply applied it to a grand hypothesis that, however plausible, had little direct experimental support." Similarly, Horace Freeland Judson records in The Eighth Day of Creation: "My mind was, that a dogma was an idea for which there was no reasonable evidence. You see?!" And Crick gave a roar of delight. "I just didn't know what dogma meant. And I could just as well have called it the 'Central Hypothesis,' or — you know. Which is what I meant to say. Dogma was just a catch phrase." Comparison with the Weismann barrier The Weismann barrier, proposed by August Weismann in 1892, distinguishes between the "immortal" germ cell lineages (the germ plasm) which produce gametes and the "disposable" somatic cells. Hereditary information moves only from germline cells to somatic cells (that is, somatic mutations are not inherited). This, before the discovery of the role or structure of DNA, does not predict the central dogma, but does anticipate its gene-centric view of life, albeit in non-molecular terms.	Biology and health sciences	Molecular biology	Biology
68300	https://en.wikipedia.org/wiki/Dominance%20%28genetics%29	Dominance (genetics)	In genetics, dominance is the phenomenon of one variant (allele) of a gene on a chromosome masking or overriding the effect of a different variant of the same gene on the other copy of the chromosome. The first variant is termed dominant and the second is called recessive. This state of having two different variants of the same gene on each chromosome is originally caused by a mutation in one of the genes, either new (de novo) or inherited. The terms autosomal dominant or autosomal recessive are used to describe gene variants on non-sex chromosomes (autosomes) and their associated traits, while those on sex chromosomes (allosomes) are termed X-linked dominant, X-linked recessive or Y-linked; these have an inheritance and presentation pattern that depends on the sex of both the parent and the child (see Sex linkage). Since there is only one copy of the Y chromosome, Y-linked traits cannot be dominant or recessive. Additionally, there are other forms of dominance, such as incomplete dominance, in which a gene variant has a partial effect compared to when it is present on both chromosomes, and co-dominance, in which different variants on each chromosome both show their associated traits. Dominance is a key concept in Mendelian inheritance and classical genetics. Letters and Punnett squares are used to demonstrate the principles of dominance in teaching, and the upper-case letters are used to denote dominant alleles and lower-case letters are used for recessive alleles. An often quoted example of dominance is the inheritance of seed shape in peas. Peas may be round, associated with allele R, or wrinkled, associated with allele r. In this case, three combinations of alleles (genotypes) are possible: RR, Rr, and rr. The RR (homozygous) individuals have round peas, and the rr (homozygous) individuals have wrinkled peas. In Rr (heterozygous) individuals, the R allele masks the presence of the r allele, so these individuals also have round peas. Thus, allele R is dominant over allele r, and allele r is recessive to allele R. Dominance is not inherent to an allele or its traits (phenotype). It is a strictly relative effect between two alleles of a given gene of any function; one allele can be dominant over a second allele of the same gene, recessive to a third, and co-dominant with a fourth. Additionally, one allele may be dominant for one trait but not others. Dominance differs from epistasis, the phenomenon of an allele of one gene masking the effect of alleles of a different gene. Background Gregor Johann Mendel, "The Father of Genetics", promulgated the idea of dominance in the 1860s. However, it was not widely known until the early twentieth century. Mendel observed that, for a variety of traits of garden peas having to do with the appearance of seeds, seed pods, and plants, there were two discrete phenotypes, such as round versus wrinkled seeds, yellow versus green seeds, red versus white flowers or tall versus short plants. When bred separately, the plants always produced the same phenotypes, generation after generation. However, when lines with different phenotypes were crossed (interbred), one and only one of the parental phenotypes showed up in the offspring (green, round, red, or tall). However, when these hybrid plants were crossed, the offspring plants showed the two original phenotypes, in a characteristic 3:1 ratio, the more common phenotype being that of the parental hybrid plants. Mendel reasoned that each parent in the first cross was a homozygote for different alleles (one parent AA and the other parent aa), that each contributed one allele to the offspring, with the result that all of these hybrids were heterozygotes (Aa), and that one of the two alleles in the hybrid cross dominated expression of the other: A masked a. The final cross between two heterozygotes (Aa X Aa) would produce AA, Aa, and aa offspring in a 1:2:1 genotype ratio with the first two classes showing the (A) phenotype, and the last showing the (a) phenotype, thereby producing the 3:1 phenotype ratio. Mendel did not use the terms gene, allele, phenotype, genotype, homozygote, and heterozygote, all of which were introduced later. He did introduce the notation of capital and lowercase letters for dominant and recessive alleles, respectively, still in use today. In 1928, British population geneticist Ronald Fisher proposed that dominance acted based on natural selection through the contribution of modifier genes. In 1929, American geneticist Sewall Wright responded by stating that dominance is simply a physiological consequence of metabolic pathways and the relative necessity of the gene involved. Types of dominance Complete dominance (Mendelian) In complete dominance, the effect of one allele in a heterozygous genotype completely masks the effect of the other. The allele that masks are considered dominant to the other allele, and the masked allele is considered recessive. When we only look at one trait determined by one pair of genes, we call it monohybrid inheritance. If the crossing is done between parents (P-generation, F0-generation) who are homozygote dominant and homozygote recessive, the offspring (F1-generation) will always have the heterozygote genotype and always present the phenotype associated with the dominant gene. However, if the F1-generation is further crossed with the F1-generation (heterozygote crossed with heterozygote) the offspring (F2-generation) will present the phenotype associated with the dominant gene ¾ times. Although heterozygote monohybrid crossing can result in two phenotype variants, it can result in three genotype variants - homozygote dominant, heterozygote and homozygote recessive, respectively. In dihybrid inheritance we look at the inheritance of two pairs of genes simultaneous. Assuming here that the two pairs of genes are located at non-homologous chromosomes, such that they are not coupled genes (see genetic linkage) but instead inherited independently. Consider now the cross between parents (P-generation) of genotypes homozygote dominant and recessive, respectively. The offspring (F1-generation) will always heterozygous and present the phenotype associated with the dominant allele variant. However, when crossing the F1-generation there are four possible phenotypic possibilities and the phenotypical ratio for the F2-generation will always be 9:3:3:1. Incomplete dominance (non-Mendelian) Incomplete dominance (also called partial dominance, semi-dominance, intermediate inheritance, or occasionally incorrectly co-dominance in reptile genetics) occurs when the phenotype of the heterozygous genotype is distinct from and often intermediate to the phenotypes of the homozygous genotypes. The phenotypic result often appears as a blended form of characteristics in the heterozygous state. For example, the snapdragon flower color is homozygous for either red or white. When the red homozygous flower is paired with the white homozygous flower, the result yields a pink snapdragon flower. The pink snapdragon is the result of incomplete dominance. A similar type of incomplete dominance is found in the four o'clock plant wherein pink color is produced when true-bred parents of white and red flowers are crossed. In quantitative genetics, where phenotypes are measured and treated numerically, if a heterozygote's phenotype is exactly between (numerically) that of the two homozygotes, the phenotype is said to exhibit no dominance at all, i.e. dominance exists only when the heterozygote's phenotype measure lies closer to one homozygote than the other. When plants of the F1 generation are self-pollinated, the phenotypic and genotypic ratio of the F2 generation will be 1:2:1 (Red:Pink:White). Co-dominance (non-Mendelian) Co-dominance occurs when the contributions of both alleles are visible in the phenotype and neither allele masks another. For example, in the ABO blood group system, chemical modifications to a glycoprotein (the H antigen) on the surfaces of blood cells are controlled by three alleles, two of which are co-dominant to each other (IA, IB) and dominant over the recessive i at the ABO locus. The IA and IB alleles produce different modifications. The enzyme coded for by IA adds an N-acetylgalactosamine to a membrane-bound H antigen. The IB enzyme adds a galactose. The i allele produces no modification. Thus the IA and IB alleles are each dominant to i (IAIA and IAi individuals both have type A blood, and IBIB and IBi individuals both have type B blood), but IAIB individuals have both modifications on their blood cells and thus have type AB blood, so the IA and IB alleles are said to be co-dominant. Another example occurs at the locus for the beta-globin component of hemoglobin, where the three molecular phenotypes of HbA/HbA, HbA/HbS, and HbS/HbS are all distinguishable by protein electrophoresis. (The medical condition produced by the heterozygous genotype is called sickle-cell trait and is a milder condition distinguishable from sickle-cell anemia, thus the alleles show incomplete dominance concerning anemia, see above). For most gene loci at the molecular level, both alleles are expressed co-dominantly, because both are transcribed into RNA. Co-dominance, where allelic products co-exist in the phenotype, is different from incomplete dominance, where the quantitative interaction of allele products produces an intermediate phenotype. For example, in co-dominance, a red homozygous flower and a white homozygous flower will produce offspring that have red and white spots. When plants of the F1 generation are self-pollinated, the phenotypic and genotypic ratio of the F2 generation will be 1:2:1 (Red:Spotted:White). These ratios are the same as those for incomplete dominance. Again, this classical terminology is inappropriate – in reality, such cases should not be said to exhibit dominance at all. Relationship to other genetic concepts Dominance can be influenced by various genetic interactions and it is essential to evaluate them when determining phenotypic outcomes. Multiple alleles, epistasis and pleiotropic genes are some factors that might influence the phenotypic outcome. Multiple alleles Although any individual of a diploid organism has at most two different alleles at a given locus, most genes exist in a large number of allelic versions in the population as a whole. This is called polymorphism, and is caused by mutations. Polymorphism can have an effect on the dominance relationship and phenotype, which is observed in the ABO blood group system. The gene responsible for human blood type have three alleles; A, B, and O, and their interactions result in different blood types based on the level of dominance the alleles expresses towards each other. Pleiotropic genes Pleiotropic genes are genes where one single gene affects two or more characters (phenotype). This means that a gene can have a dominant effect on one trait, but a more recessive effect on another trait. Epistasis Epistasis is interactions between multiple alleles at different loci. Easily said, several genes for one phenotype. The dominance relationship between alleles involved in epistatic interactions can influence the observed phenotypic ratios in offspring.	Biology and health sciences	Genetics	Biology
68316	https://en.wikipedia.org/wiki/Heat%20pump	Heat pump	A heat pump is a device that uses energy (usually electricity) to transfer heat from a colder place to a warmer place. Specifically, the heat pump transfers thermal energy using a heat pump and refrigeration cycle, cooling the cool space and warming the warm space. In winter a heat pump can move heat from the cool outdoors to warm a house; the pump may also be designed to move heat from the house to the warmer outdoors in summer. As they transfer heat rather than generating heat, they are more energy-efficient than heating by gas boiler, and also good for cooling a home. A gaseous refrigerant is compressed so its pressure and temperature rise. When operating as a heater in cold weather, the warmed gas flows to a heat exchanger in the indoor space where some of its thermal energy is transferred to that indoor space, causing the gas to condense to its liquid state. The liquified refrigerant flows to a heat exchanger in the outdoor space where the pressure falls, the liquid evaporates and the temperature of the gas falls. It is now colder than the temperature of the outdoor space being used as a heat source. It can again take up energy from the heat source, be compressed and repeat the cycle. Air source heat pumps are the most common models, while other types include ground source heat pumps, water source heat pumps and exhaust air heat pumps. Large-scale heat pumps are also used in district heating systems. The efficiency of a heat pump is expressed as a coefficient of performance (COP), or seasonal coefficient of performance (SCOP). The higher the number, the more efficient a heat pump is. For example, an air-to-water heat pump that produces 6kW at a SCOP of 4.62 will give over 4kW of energy into a heating system for every kilowatt of energy that the heat pump uses itself to operate. When used for space heating, heat pumps are typically more energy-efficient than electric resistance and other heaters. Because of their high efficiency and the increasing share of fossil-free sources in electrical grids, heat pumps are playing a role in climate change mitigation. Consuming 1 kWh of electricity, they can transfer 1 to 4.5 kWh of thermal energy into a building. The carbon footprint of heat pumps depends on how electricity is generated, but they usually reduce emissions. Heat pumps could satisfy over 80% of global space and water heating needs with a lower carbon footprint than gas-fired condensing boilers: however, in 2021 they only met 10%. Principle of operation Heat flows spontaneously from a region of higher temperature to a region of lower temperature. Heat does not flow spontaneously from lower temperature to higher, but it can be made to flow in this direction if work is performed. The work required to transfer a given amount of heat is usually much less than the amount of heat; this is the motivation for using heat pumps in applications such as the heating of water and the interior of buildings. The amount of work required to drive an amount of heat Q from a lower-temperature reservoir such as ambient air to a higher-temperature reservoir such as the interior of a building is: where is the work performed on the working fluid by the heat pump's compressor. is the heat transferred from the lower-temperature reservoir to the higher-temperature reservoir. is the instantaneous coefficient of performance for the heat pump at the temperatures prevailing in the reservoirs at one instant. The coefficient of performance of a heat pump is greater than one so the work required is less than the heat transferred, making a heat pump a more efficient form of heating than electrical resistance heating. As the temperature of the higher-temperature reservoir increases in response to the heat flowing into it, the coefficient of performance decreases, causing an increasing amount of work to be required for each unit of heat being transferred. The coefficient of performance, and the work required by a heat pump can be calculated easily by considering an ideal heat pump operating on the reversed Carnot cycle: If the low-temperature reservoir is at a temperature of and the interior of the building is at the relevant coefficient of performance is 27. This means only 1 joule of work is required to transfer 27 joules of heat from a reservoir at 270 K to another at 280 K. The one joule of work ultimately ends up as thermal energy in the interior of the building so for each 27 joules of heat that are removed from the low-temperature reservoir, 28 joules of heat are added to the building interior, making the heat pump even more attractive from an efficiency perspective. As the temperature of the interior of the building rises progressively to the coefficient of performance falls progressively to 9. This means each joule of work is responsible for transferring 9 joules of heat out of the low-temperature reservoir and into the building. Again, the 1 joule of work ultimately ends up as thermal energy in the interior of the building so 10 joules of heat are added to the building interior. This is the theoretical amount of heat pumped but in practice it will be less for various reasons, for example if the outside unit has been installed where there is not enough airflow. More data sharing with owners and academics—perhaps from heat meters—could improve efficiency in the long run. History Milestones: 1748 William Cullen demonstrates artificial refrigeration. 1834 Jacob Perkins patents a design for a practical refrigerator using dimethyl ether. 1852 Lord Kelvin describes the theory underlying heat pumps. 1855–1857 Peter von Rittinger develops and builds the first heat pump. 1877 In the period before 1875, heat pumps were for the time being pursued for vapour compression evaporation (open heat pump process) in salt works with their obvious advantages for saving wood and coal. In 1857, Peter von Rittinger was the first to try to implement the idea of vapor compression in a small pilot plant. Presumably inspired by Rittinger's experiments in Ebensee, Antoine-Paul Piccard from the University of Lausanne and the engineer J. H. Weibel from the Weibel–Briquet company in Geneva built the world's first really functioning vapor compression system with a two-stage piston compressor. In 1877 this first heat pump in Switzerland was installed in the Bex salt works. 1928 Aurel Stodola constructs a closed-loop heat pump (water source from Lake Geneva) which provides heating for the Geneva city hall to this day. 1937–1945 During the First World War, fuel prices were very high in Switzerland but it had plenty of hydropower. In the period before and especially during the Second World War, when neutral Switzerland was completely surrounded by fascist-ruled countries, the coal shortage became alarming again. Thanks to their leading position in energy technology, the Swiss companies Sulzer, Escher Wyss and Brown Boveri built and put in operation around 35 heat pumps between 1937 and 1945. The main heat sources were lake water, river water, groundwater, and waste heat. Particularly noteworthy are the six historic heat pumps from the city of Zurich with heat outputs from 100 kW to 6 MW. An international milestone is the heat pump built by Escher Wyss in 1937/38 to replace the wood stoves in the City Hall of Zurich. To avoid noise and vibrations, a recently developed rotary piston compressor was used. This historic heat pump heated the town hall for 63 years until 2001. Only then was it replaced by a new, more efficient heat pump. 1945 John Sumner, City Electrical Engineer for Norwich, installs an experimental water-source heat pump fed central heating system, using a nearby river to heat new Council administrative buildings. It had a seasonal efficiency ratio of 3.42, average thermal delivery of 147 kW, and peak output of 234 kW. 1948 Robert C. Webber is credited as developing and building the first ground-source heat pump. 1951 First large scale installation—the Royal Festival Hall in London is opened with a town gas-powered reversible water-source heat pump, fed by the Thames, for both winter heating and summer cooling needs. 2019 The Kigali Amendment to phase out harmful refrigerants takes effect. Types Air-source Ground source Heat recovery ventilation Exhaust air heat pumps extract heat from the exhaust air of a building and require mechanical ventilation. Two classes exist: Exhaust air-air heat pumps transfer heat to intake air. Exhaust air-water heat pumps transfer heat to a heating circuit that includes a tank of domestic hot water. Solar-assisted Water-source A water-source heat pump works in a similar manner to a ground-source heat pump, except that it takes heat from a body of water rather than the ground. The body of water does, however, need to be large enough to be able to withstand the cooling effect of the unit without freezing or creating an adverse effect for wildlife. The largest water-source heat pump was installed in the Danish town of Esbjerg in 2023. Others A thermoacoustic heat pump operates as a thermoacoustic heat engine without refrigerant but instead uses a standing wave in a sealed chamber driven by a loudspeaker to achieve a temperature difference across the chamber. Electrocaloric heat pumps are solid state. Applications The International Energy Agency estimated that, as of 2021, heat pumps installed in buildings have a combined capacity of more than 1000 GW. They are used for heating, ventilation, and air conditioning (HVAC) and may also provide domestic hot water and tumble clothes drying. The purchase costs are supported in various countries by consumer rebates. Space heating and sometimes also cooling In HVAC applications, a heat pump is typically a vapor-compression refrigeration device that includes a reversing valve and optimized heat exchangers so that the direction of heat flow (thermal energy movement) may be reversed. The reversing valve switches the direction of refrigerant through the cycle and therefore the heat pump may deliver either heating or cooling to a building. Because the two heat exchangers, the condenser and evaporator, must swap functions, they are optimized to perform adequately in both modes. Therefore, the Seasonal Energy Efficiency Rating (SEER in the US) or European seasonal energy efficiency ratio of a reversible heat pump is typically slightly less than those of two separately optimized machines. For equipment to receive the US Energy Star rating, it must have a rating of at least 14 SEER. Pumps with ratings of 18 SEER or above are considered highly efficient. The highest efficiency heat pumps manufactured are up to 24 SEER. Heating seasonal performance factor (in the US) or Seasonal Performance Factor (in Europe) are ratings of heating performance. The SPF is Total heat output per annum / Total electricity consumed per annum in other words the average heating COP over the year. Window mounted heat pump Window mounted heat pumps run on standard 120v AC outlets and provide heating, cooling, and humidity control. They are more efficient with lower noise levels, condensation management, and a smaller footprint than window mounted air conditioners that just do cooling. Water heating In water heating applications, heat pumps may be used to heat or preheat water for swimming pools, homes or industry. Usually heat is extracted from outdoor air and transferred to an indoor water tank. District heating Large (megawatt-scale) heat pumps are used for district heating. However about 90% of district heat is from fossil fuels. In Europe, heat pumps account for a mere 1% of heat supply in district heating networks but several countries have targets to decarbonise their networks between 2030 and 2040. Possible sources of heat for such applications are sewage water, ambient water (e.g. sea, lake and river water), industrial waste heat, geothermal energy, flue gas, waste heat from district cooling and heat from solar seasonal thermal energy storage. Large-scale heat pumps for district heating combined with thermal energy storage offer high flexibility for the integration of variable renewable energy. Therefore, they are regarded as a key technology for limiting climate change by phasing out fossil fuels. They are also a crucial element of systems which can both heat and cool districts. Industrial heating There is great potential to reduce the energy consumption and related greenhouse gas emissions in industry by application of industrial heat pumps, for example for process heat. Short payback periods of less than 2 years are possible, while achieving a high reduction of emissions (in some cases more than 50%). Industrial heat pumps can heat up to 200 °C, and can meet the heating demands of many light industries. In Europe alone, 15 GW of heat pumps could be installed in 3,000 facilities in the paper, food and chemicals industries. Performance The performance of a heat pump is determined by the ability of the pump to extract heat from a low temperature environment (the source) and deliver it to a higher temperature environment (the sink). Performance varies, depending on installation details, temperature differences, site elevation, location on site, pipe runs, flow rates, and maintenance. In general, heat pumps work most efficiently (that is, the heat output produced for a given energy input) when the difference between the heat source and the heat sink is small. When using a heat pump for space or water heating, therefore, the heat pump will be most efficient in mild conditions, and decline in efficiency on very cold days. Performance metrics supplied to consumers attempt to take this variation into account. Common performance metrics are the SEER (in cooling mode) and seasonal coefficient of performance (SCOP) (commonly used just for heating), although SCOP can be used for both modes of operation. Larger values of either metric indicate better performance. When comparing the performance of heat pumps, the term performance is preferred to efficiency, with coefficient of performance (COP) being used to describe the ratio of useful heat movement per work input. An electrical resistance heater has a COP of 1.0, which is considerably lower than a well-designed heat pump which will typically have a COP of 3 to 5 with an external temperature of 10 °C and an internal temperature of 20 °C. Because the ground is a constant temperature source, a ground-source heat pump is not subjected to large temperature fluctuations, and therefore is the most energy-efficient type of heat pump. The "seasonal coefficient of performance" (SCOP) is a measure of the aggregate energy efficiency measure over a period of one year which is dependent on regional climate. One framework for this calculation is given by the Commission Regulation (EU) No. 813/2013. A heat pump's operating performance in cooling mode is characterized in the US by either its energy efficiency ratio (EER) or seasonal energy efficiency ratio (SEER), both of which have units of BTU/(h·W) (note that 1 BTU/(h·W) = 0.293 W/W) and larger values indicate better performance. Carbon footprint The carbon footprint of heat pumps depends on their individual efficiency and how electricity is produced. An increasing share of low-carbon energy sources such as wind and solar will lower the impact on the climate. In most settings, heat pumps will reduce emissions compared to heating systems powered by fossil fuels. In regions accounting for 70% of world energy consumption, the emissions savings of heat pumps compared with a high-efficiency gas boiler are on average above 45% and reach 80% in countries with cleaner electricity mixes. These values can be improved by 10 percentage points, respectively, with alternative refrigerants. In the United States, 70% of houses could reduce emissions by installing a heat pump. The rising share of renewable electricity generation in many countries is set to increase the emissions savings from heat pumps over time. Heating systems powered by green hydrogen are also low-carbon and may become competitors, but are much less efficient due to the energy loss associated with hydrogen conversion, transport and use. In addition, not enough green hydrogen is expected to be available before the 2030s or 2040s. Operation Vapor-compression uses a circulating refrigerant as the medium which absorbs heat from one space, compresses it thereby increasing its temperature before releasing it in another space. The system normally has eight main components: a compressor, a reservoir, a reversing valve which selects between heating and cooling mode, two thermal expansion valves (one used when in heating mode and the other when used in cooling mode) and two heat exchangers, one associated with the external heat source/sink and the other with the interior. In heating mode the external heat exchanger is the evaporator and the internal one being the condenser; in cooling mode the roles are reversed. Circulating refrigerant enters the compressor in the thermodynamic state known as a saturated vapor and is compressed to a higher pressure, resulting in a higher temperature as well. The hot, compressed vapor is then in the thermodynamic state known as a superheated vapor and it is at a temperature and pressure at which it can be condensed with either cooling water or cooling air flowing across the coil or tubes. In heating mode this heat is used to heat the building using the internal heat exchanger, and in cooling mode this heat is rejected via the external heat exchanger. The condensed, liquid refrigerant, in the thermodynamic state known as a saturated liquid, is next routed through an expansion valve where it undergoes an abrupt reduction in pressure. That pressure reduction results in the adiabatic flash evaporation of a part of the liquid refrigerant. The auto-refrigeration effect of the adiabatic flash evaporation lowers the temperature of the liquid and-vapor refrigerant mixture to where it is colder than the temperature of the enclosed space to be refrigerated. The cold mixture is then routed through the coil or tubes in the evaporator. A fan circulates the warm air in the enclosed space across the coil or tubes carrying the cold refrigerant liquid and vapor mixture. That warm air evaporates the liquid part of the cold refrigerant mixture. At the same time, the circulating air is cooled and thus lowers the temperature of the enclosed space to the desired temperature. The evaporator is where the circulating refrigerant absorbs and removes heat which is subsequently rejected in the condenser and transferred elsewhere by the water or air used in the condenser. To complete the refrigeration cycle, the refrigerant vapor from the evaporator is again a saturated vapor and is routed back into the compressor. Over time, the evaporator may collect ice or water from ambient humidity. The ice is melted through defrosting cycle. An internal heat exchanger is either used to heat/cool the interior air directly or to heat water that is then circulated through radiators or underfloor heating circuit to either heat or cool the buildings. Improvement of coefficient of performance by subcooling Heat input can be improved if the refrigerant enters the evaporator with a lower vapor content. This can be achieved by cooling the liquid refrigerant after condensation. The gaseous refrigerant condenses on the heat exchange surface of the condenser. To achieve a heat flow from the gaseous flow center to the wall of the condenser, the temperature of the liquid refrigerant must be lower than the condensation temperature. Additional subcooling can be achieved by heat exchange between relatively warm liquid refrigerant leaving the condenser and the cooler refrigerant vapor emerging from the evaporator. The enthalpy difference required for the subcooling leads to the superheating of the vapor drawn into the compressor. When the increase in cooling achieved by subcooling is greater that the compressor drive input required to overcome the additional pressure losses, such a heat exchange improves the coefficient of performance. One disadvantage of the subcooling of liquids is that the difference between the condensing temperature and the heat-sink temperature must be larger. This leads to a moderately high pressure difference between condensing and evaporating pressure, whereby the compressor energy increases. Refrigerant choice Pure refrigerants can be divided into organic substances (hydrocarbons (HCs), chlorofluorocarbons (CFCs), hydrochlorofluorocarbons (HCFCs), hydrofluorocarbons (HFCs), hydrofluoroolefins (HFOs), and HCFOs), and inorganic substances (ammonia (), carbon dioxide (), and water ()). Their boiling points are usually below −25 °C. In the past 200 years, the standards and requirements for new refrigerants have changed. Nowadays low global warming potential (GWP) is required, in addition to all the previous requirements for safety, practicality, material compatibility, appropriate atmospheric life, and compatibility with high-efficiency products. By 2022, devices using refrigerants with a very low GWP still have a small market share but are expected to play an increasing role due to enforced regulations, as most countries have now ratified the Kigali Amendment to ban HFCs. Isobutane (R600A) and propane (R290) are far less harmful to the environment than conventional hydrofluorocarbons (HFC) and are already being used in air-source heat pumps. Propane may be the most suitable for high temperature heat pumps. Ammonia (R717) and carbon dioxide (R-744) also have a low GWP. smaller heat pumps are not widely available and research and development of them continues. A 2024 report said that refrigerants with GWP are vulnerable to further international restrictions. Until the 1990s, heat pumps, along with fridges and other related products used chlorofluorocarbons (CFCs) as refrigerants, which caused major damage to the ozone layer when released into the atmosphere. Use of these chemicals was banned or severely restricted by the Montreal Protocol of August 1987. Replacements, including R-134a and R-410A, are hydrofluorocarbons (HFC) with similar thermodynamic properties with insignificant ozone depletion potential (ODP) but had problematic GWP. HFCs are powerful greenhouse gases which contribute to climate change. Dimethyl ether (DME) also gained in popularity as a refrigerant in combination with R404a. More recent refrigerants include difluoromethane (R32) with a lower GWP, but still over 600. Devices with R-290 refrigerant (propane) are expected to play a key role in the future. The 100-year GWP of propane, at 0.02, is extremely low and is approximately 7000 times less than R-32. However, the flammability of propane requires additional safety measures: the maximum safe charges have been set significantly lower than for lower flammability refrigerants (only allowing approximately 13.5 times less refrigerant in the system than R-32). This means that R-290 is not suitable for all situations or locations. Nonetheless, by 2022, an increasing number of devices with R-290 were offered for domestic use, especially in Europe. At the same time, HFC refrigerants still dominate the market. Recent government mandates have seen the phase-out of R-22 refrigerant. Replacements such as R-32 and R-410A are being promoted as environmentally friendly but still have a high GWP. A heat pump typically uses 3 kg of refrigerant. With R-32 this amount still has a 20-year impact equivalent to 7 tons of , which corresponds to two years of natural gas heating in an average household. Refrigerants with a high ODP have already been phased out. Government incentives Financial incentives aim to protect consumers from high fossil gas costs and to reduce greenhouse gas emissions, and are currently available in more than 30 countries around the world, covering more than 70% of global heating demand in 2021. Australia Food processors, brewers, petfood producers and other industrial energy users are exploring whether it is feasible to use renewable energy to produce industrial-grade heat. Process heating accounts for the largest share of onsite energy use in Australian manufacturing, with lower-temperature operations like food production particularly well-suited to transition to renewables. To help producers understand how they could benefit from making the switch, the Australian Renewable Energy Agency (ARENA) provided funding to the Australian Alliance for Energy Productivity (A2EP) to undertake pre-feasibility studies at a range of sites around Australia, with the most promising locations advancing to full feasibility studies. In an effort to incentivize energy efficiency and reduce environmental impact, the Australian states of Victoria, New South Wales, and Queensland have implemented rebate programs targeting the upgrade of existing hot water systems. These programs specifically encourage the transition from traditional gas or electric systems to heat pump based systems. Canada In 2022, the Canada Greener Homes Grant provides up to $5000 for upgrades (including certain heat pumps), and $600 for energy efficiency evaluations. China Purchase subsidies in rural areas in the 2010s reduced burning coal for heating, which had been causing ill health. In the 2024 report by the International Energy Agency (IEA) titled "The Future of Heat Pumps in China," it is highlighted that China, as the world's largest market for heat pumps in buildings, plays a critical role in the global industry. The country accounts for over one-quarter of global sales, with a 12% increase in 2023 alone, despite a global sales dip of 3% the same year. Heat pumps are now used in approximately 8% of all heating equipment sales for buildings in China as of 2022, and they are increasingly becoming the norm in central and southern regions for both heating and cooling. Despite their higher upfront costs and relatively low awareness, heat pumps are favored for their energy efficiency, consuming three to five times less energy than electric heaters or fossil fuel-based solutions. Currently, decentralized heat pumps installed in Chinese buildings represent a quarter of the global installed capacity, with a total capacity exceeding 250 GW, which covers around 4% of the heating needs in buildings. Under the Announced Pledges Scenario (APS), which aligns with China's carbon neutrality goals, the capacity is expected to reach 1,400 GW by 2050, meeting 25% of heating needs. This scenario would require an installation of about 100 GW of heat pumps annually until 2050. Furthermore, the heat pump sector in China employs over 300,000 people, with employment numbers expected to double by 2050, underscoring the importance of vocational training for industry growth. This robust development in the heat pump market is set to play a significant role in reducing direct emissions in buildings by 30% and cutting PM2.5 emissions from residential heating by nearly 80% by 2030. European Union To speed up the deployment rate of heat pumps, the European Commission launched the Heat Pump Accelerator Platform in November 2024. It will encourage industry experts, policymakers, and stakeholders to collaborate, share best practices and ideas, and jointly discuss measures that promote sustainable heating solutions. United Kingdom Until 2027 fixed heat pumps have no Value Added Tax (VAT). the installation cost of a heat pump is more than a gas boiler, but with the "Boiler Upgrade Scheme" government grant and assuming electricity/gas costs remain similar their lifetime costs would be similar on average. However lifetime cost relative to a gas boiler varies considerably depending on several factors, such as the quality of the heat pump installation and the tariff used. In 2024 England was criticised for still allowing new homes to be built with gas boilers, unlike some other counties where this is banned. United States The High-efficiency Electric Home Rebate Program was created in 2022 to award grants to State energy offices and Indian Tribes in order to establish state-wide high-efficiency electric-home rebates. Effective immediately, American households are eligible for a tax credit to cover the costs of buying and installing a heat pump, up to $2,000. Starting in 2023, low- and moderate-level income households will be eligible for a heat-pump rebate of up to $8,000. In 2022, more heat pumps were sold in the United States than natural gas furnaces. In November 2023 Biden's administration allocated 169 million dollars from the Inflation Reduction Act to speed up production of heat pumps. It used the Defense Production Act to do so, because according to the administration, energy that is better for the climate is also better for national security.	Technology	Heating and cooling	null
68326	https://en.wikipedia.org/wiki/Extended%20periodic%20table	Extended periodic table	An extended periodic table theorizes about chemical elements beyond those currently known and proven. The element with the highest atomic number known is oganesson (Z = 118), which completes the seventh period (row) in the periodic table. All elements in the eighth period and beyond thus remain purely hypothetical. Elements beyond 118 will be placed in additional periods when discovered, laid out (as with the existing periods) to illustrate periodically recurring trends in the properties of the elements. Any additional periods are expected to contain more elements than the seventh period, as they are calculated to have an additional so-called g-block, containing at least 18 elements with partially filled g-orbitals in each period. An eight-period table containing this block was suggested by Glenn T. Seaborg in 1969. The first element of the g-block may have atomic number 121, and thus would have the systematic name unbiunium. Despite many searches, no elements in this region have been synthesized or discovered in nature. According to the orbital approximation in quantum mechanical descriptions of atomic structure, the g-block would correspond to elements with partially filled g-orbitals, but spin–orbit coupling effects reduce the validity of the orbital approximation substantially for elements of high atomic number. Seaborg's version of the extended period had the heavier elements following the pattern set by lighter elements, as it did not take into account relativistic effects. Models that take relativistic effects into account predict that the pattern will be broken. Pekka Pyykkö and Burkhard Fricke used computer modeling to calculate the positions of elements up to Z = 172, and found that several were displaced from the Madelung rule. As a result of uncertainty and variability in predictions of chemical and physical properties of elements beyond 120, there is currently no consensus on their placement in the extended periodic table. Elements in this region are likely to be highly unstable with respect to radioactive decay and undergo alpha decay or spontaneous fission with extremely short half-lives, though element 126 is hypothesized to be within an island of stability that is resistant to fission but not to alpha decay. Other islands of stability beyond the known elements may also be possible, including one theorised around element 164, though the extent of stabilizing effects from closed nuclear shells is uncertain. It is not clear how many elements beyond the expected island of stability are physically possible, whether period 8 is complete, or if there is a period 9. The International Union of Pure and Applied Chemistry (IUPAC) defines an element to exist if its lifetime is longer than 10−14 seconds (0.01 picoseconds, or 10 femtoseconds), which is the time it takes for the nucleus to form an electron cloud. As early as 1940, it was noted that a simplistic interpretation of the relativistic Dirac equation runs into problems with electron orbitals at Z > 1/α ≈ 137, suggesting that neutral atoms cannot exist beyond element 137, and that a periodic table of elements based on electron orbitals therefore breaks down at this point. On the other hand, a more rigorous analysis calculates the analogous limit to be Z ≈ 168–172 where the 1s subshell dives into the Dirac sea, and that it is instead not neutral atoms that cannot exist beyond this point, but bare nuclei, thus posing no obstacle to the further extension of the periodic system. Atoms beyond this critical atomic number are called supercritical atoms. History Elements beyond the actinides were first proposed to exist as early as 1895, when Danish chemist Hans Peter Jørgen Julius Thomsen predicted that thorium and uranium formed part of a 32-element period which would end at a chemically inactive element with atomic weight 292 (not far from the 294 for the only known isotope of oganesson). In 1913, Swedish physicist Johannes Rydberg similarly predicted that the next noble gas after radon would have atomic number 118, and purely formally derived even heavier congeners of radon at Z = 168, 218, 290, 362, and 460, exactly where the Aufbau principle would predict them to be. In 1922, Niels Bohr predicted the electronic structure of this next noble gas at Z = 118, and suggested that the reason why elements beyond uranium were not seen in nature was because they were too unstable. The German physicist and engineer Richard Swinne published a review paper in 1926 containing predictions on the transuranic elements (he may have coined the term) in which he anticipated modern predictions of an island of stability: he first hypothesised in 1914 that half-lives should not decrease strictly with atomic number, but suggested instead that there might be some longer-lived elements at Z = 98–102 and Z = 108–110, and speculated that such elements might exist in the Earth's core, in iron meteorites, or in the ice caps of Greenland where they had been locked up from their supposed cosmic origin. By 1955, these elements were called superheavy elements. The first predictions on properties of undiscovered superheavy elements were made in 1957, when the concept of nuclear shells was first explored and an island of stability was theorized to exist around element 126. In 1967, more rigorous calculations were performed, and the island of stability was theorized to be centered at the then-undiscovered flerovium (element 114); this and other subsequent studies motivated many researchers to search for superheavy elements in nature or attempt to synthesize them at accelerators. Many searches for superheavy elements were conducted in the 1970s, all with negative results. , synthesis has been attempted for every element up to and including unbiseptium (Z = 127), except unbitrium (Z = 123), with the heaviest successfully synthesized element being oganesson in 2002 and the most recent discovery being that of tennessine in 2010. As some superheavy elements were predicted to lie beyond the seven-period periodic table, an additional eighth period containing these elements was first proposed by Glenn T. Seaborg in 1969. This model continued the pattern in established elements and introduced a new g-block and superactinide series beginning at element 121, raising the number of elements in period 8 compared to known periods. These early calculations failed to consider relativistic effects that break down periodic trends and render simple extrapolation impossible, however. In 1971, Fricke calculated the periodic table up to Z = 172, and discovered that some elements indeed had different properties that break the established pattern, and a 2010 calculation by Pekka Pyykkö also noted that several elements might behave differently than expected. It is unknown how far the periodic table might extend beyond the known 118 elements, as heavier elements are predicted to be increasingly unstable. Glenn T. Seaborg suggested that practically speaking, the end of the periodic table might come as early as around Z = 120 due to nuclear instability. Predicted structures of an extended periodic table There is currently no consensus on the placement of elements beyond atomic number 120 in the periodic table. All hypothetical elements are given an International Union of Pure and Applied Chemistry (IUPAC) systematic element name, for use until the element has been discovered, confirmed, and an official name is approved. These names are typically not used in the literature, and the elements are instead referred to by their atomic numbers; hence, element 164 is usually not called "unhexquadium" or "Uhq" (the systematic name and symbol), but rather "element 164" with symbol "164", "(164)", or "E164". Aufbau principle At element 118, the orbitals 1s, 2s, 2p, 3s, 3p, 3d, 4s, 4p, 4d, 4f, 5s, 5p, 5d, 5f, 6s, 6p, 6d, 7s and 7p are assumed to be filled, with the remaining orbitals unfilled. A simple extrapolation from the Aufbau principle would predict the eighth row to fill orbitals in the order 8s, 5g, 6f, 7d, 8p; but after element 120, the proximity of the electron shells makes placement in a simple table problematic. Fricke Not all models show the higher elements following the pattern established by lighter elements. Burkhard Fricke et al., who carried out calculations up to element 184 in an article published in 1971, also found some elements to be displaced from the Madelung energy-ordering rule as a result of overlapping orbitals; this is caused by the increasing role of relativistic effects in heavy elements (They describe chemical properties up to element 184, but only draw a table to element 172.) Fricke et al.'s format is more focused on formal electron configurations than likely chemical behaviour. They place elements 156–164 in groups 4–12 because formally their configurations should be 7d2 through 7d10. However, they differ from the previous d-elements in that the 8s shell is not available for chemical bonding: instead, the 9s shell is. Thus element 164 with 7d109s0 is noted by Fricke et al. to be analogous to palladium with 4d105s0, and they consider elements 157–172 to have chemical analogies to groups 3–18 (though they are ambivalent on whether elements 165 and 166 are more like group 1 and 2 elements or more like group 11 and 12 elements, respectively). Thus, elements 157–164 are placed in their table in a group that the authors do not think is chemically most analogous. Nefedov , Trzhaskovskaya, and Yarzhemskii carried out calculations up to 164 (results published in 2006). They considered elements 158 through 164 to be homologues of groups 4 through 10, and not 6 through 12, noting similarities of electron configurations to the period 5 transition metals (e.g. element 159 7d49s1 vs Nb 4d45s1, element 160 7d59s1 vs Mo 4d55s1, element 162 7d79s1 vs Ru 4d75s1, element 163 7d89s1 vs Rh 4d85s1, element 164 7d109s0 vs Pd 4d105s0). They thus agree with Fricke et al. on the chemically most analogous groups, but differ from them in that Nefedov et al. actually place elements in the chemically most analogous groups. Rg and Cn are given an asterisk to reflect differing configurations from Au and Hg (in the original publication they are drawn as being displaced in the third dimension). In fact Cn probably has an analogous configuration to Hg, and the difference in configuration between Pt and Ds is not marked. Pyykkö Pekka Pyykkö used computer modeling to calculate the positions of elements up to Z = 172 and their possible chemical properties in an article published in 2011. He reproduced the orbital order of Fricke et al., and proposed a refinement of their table by formally assigning slots to elements 121–164 based on ionic configurations. In order to bookkeep the electrons, Pyykkö places some elements out of order: thus 139 and 140 are placed in groups 13 and 14 to reflect that the 8p1/2 shell needs to fill, and he distinguishes separate , 8p1/2, and 6f series. Fricke et al. and Nefedov et al. do not attempt to break up these series. Kulsha Computational chemist Andrey Kulsha has suggested two forms of the extended periodic table up to 172 that build on and refine Nefedov et al.'s versions up to 164 with reference to Pyykkö's calculations. Based on their likely chemical properties, elements 157–172 are placed by both forms as eighth-period congeners of yttrium through xenon in the fifth period; this extends Nefedov et al.'s placement of 157–164 under yttrium through palladium, and agrees with the chemical analogies given by Fricke et al. Kulsha suggested two ways to deal with elements 121–156, that lack precise analogues among earlier elements. In his first form (2011, after Pyykkö's paper was published), elements 121–138 and 139–156 are placed as two separate rows (together called "ultransition elements"), related by the addition of a 5g18 subshell into the core, as according to Pyykkö's calculations of oxidation states, they should, respectively, mimic lanthanides and actinides. In his second suggestion (2016), elements 121–142 form a g-block (as they have 5g activity), while elements 143–156 form an f-block placed under actinium through nobelium. Thus, period 8 emerges with 54 elements, and the next noble element after 118 is 172. Smits et al. In 2023 Smits, Düllmann, Indelicato, Nazarewicz, and Schwerdtfeger made another attempt to place elements from 119 to 170 in the periodic table based on their electron configurations. The configurations of a few elements (121–124 and 168) did not allow them to be placed unambiguously. Element 145 appears twice, some places have double occupancy, and others are empty. Searches for undiscovered elements Synthesis attempts Attempts have been made to synthesise the period 8 elements up to unbiseptium, except unbitrium. All such attempts have been unsuccessful. An attempt to synthesise ununennium, the first period 8 element, is ongoing . Ununennium (E119) The synthesis of element 119 (ununennium) was first attempted in 1985 by bombarding a target of einsteinium-254 with calcium-48 ions at the superHILAC accelerator at Berkeley, California: + → 119* → no atoms No atoms were identified, leading to a limiting cross section of 300 nb. Later calculations suggest that the cross section of the 3n reaction (which would result in 119 and three neutrons as products) would actually be six hundred thousand times lower than this upper bound, at 0.5 pb. From April to September 2012, an attempt to synthesize the isotopes 119 and 119 was made by bombarding a target of berkelium-249 with titanium-50 at the GSI Helmholtz Centre for Heavy Ion Research in Darmstadt, Germany. Based on the theoretically predicted cross section, it was expected that an ununennium atom would be synthesized within five months of the beginning of the experiment. Moreover, as berkelium-249 decays to californium-249 (the next element) with a short half-life of 327 days, this allowed elements 119 and 120 to be searched for simultaneously. + → 119* → no atoms The experiment was originally planned to continue to November 2012, but was stopped early to make use of the Bk target to confirm the synthesis of tennessine (thus changing the projectiles to Ca). This reaction of Bk + Ti was predicted to be the most favorable practical reaction for formation of element 119, as it is rather asymmetrical, though also somewhat cold. (Es + Ca would be superior, but preparing milligram quantities of Es for a target is difficult.) Nevertheless, the necessary change from the "silver bullet" Ca to Ti divides the expected yield of element 119 by about twenty, as the yield is strongly dependent on the asymmetry of the fusion reaction. Due to the predicted short half-lives, the GSI team used new "fast" electronics capable of registering decay events within microseconds. No atoms of element 119 were identified, implying a limiting cross section of 70 fb. The predicted actual cross section is around 40 fb, which is at the limits of current technology. The team at RIKEN in Wakō, Japan began bombarding curium-248 targets with a vanadium-51 beam in January 2018 to search for element 119. Curium was chosen as a target, rather than heavier berkelium or californium, as these heavier targets are difficult to prepare. The Cm targets were provided by Oak Ridge National Laboratory. RIKEN developed a high-intensity vanadium beam. The experiment began at a cyclotron while RIKEN upgraded its linear accelerators; the upgrade was completed in 2020. Bombardment may be continued with both machines until the first event is observed; the experiment is currently running intermittently for at least 100 days a year. The RIKEN team's efforts are being financed by the Emperor of Japan. The team at the JINR plans to attempt synthesis of element 119 in the future, probably via the Am + Cr reaction, but a precise timeframe has not been publicly released. Unbinilium (E120) Following their success in obtaining oganesson by the reaction between 249Cf and 48Ca in 2006, the team at the Joint Institute for Nuclear Research (JINR) in Dubna started similar experiments in March–April 2007, in hope of creating element 120 (unbinilium) from nuclei of 58Fe and 244Pu. Isotopes of unbinilium are predicted to have alpha decay half-lives of the order of microseconds. Initial analysis revealed that no atoms of element 120 were produced, providing a limit of 400 fb for the cross section at the energy studied. + → 302120* → no atoms The Russian team planned to upgrade their facilities before attempting the reaction again. In April 2007, the team at the GSI Helmholtz Centre for Heavy Ion Research in Darmstadt, Germany, attempted to create element 120 using uranium-238 and nickel-64: + → 302120* → no atoms No atoms were detected, providing a limit of 1.6 pb for the cross section at the energy provided. The GSI repeated the experiment with higher sensitivity in three separate runs in April–May 2007, January–March 2008, and September–October 2008, all with negative results, reaching a cross section limit of 90 fb. In June–July 2010, and again in 2011, after upgrading their equipment to allow the use of more radioactive targets, scientists at the GSI attempted the more asymmetrical fusion reaction: + → 302120 → no atoms It was expected that the change in reaction would quintuple the probability of synthesizing element 120, as the yield of such reactions is strongly dependent on their asymmetry. Three correlated signals were observed that matched the predicted alpha decay energies of 299120 and its daughter 295Og, as well as the experimentally known decay energy of its granddaughter 291Lv. However, the lifetimes of these possible decays were much longer than expected, and the results could not be confirmed. In August–October 2011, a different team at the GSI using the TASCA facility tried a new, even more asymmetrical reaction: + → 299120* → no atoms This was also tried unsuccessfully the next year during the aforementioned attempt to make element 119 in the 249Bk+50Ti reaction, as 249Bk decays to 249Cf. Because of its asymmetry, the reaction between 249Cf and 50Ti was predicted to be the most favorable practical reaction for synthesizing unbinilium, although it is also somewhat cold. No unbinilium atoms were identified, implying a limiting cross-section of 200 fb. Jens Volker Kratz predicted the actual maximum cross-section for producing element 120 by any of these reactions to be around 0.1 fb; in comparison, the world record for the smallest cross section of a successful reaction was 30 fb for the reaction 209Bi(70Zn,n)278Nh, and Kratz predicted a maximum cross-section of 20 fb for producing the neighbouring element 119. If these predictions are accurate, then synthesizing element 119 would be at the limits of current technology, and synthesizing element 120 would require new methods. In May 2021, the JINR announced plans to investigate the 249Cf+50Ti reaction in their new facility. However, the 249Cf target would have had to be made by the Oak Ridge National Laboratory in the United States, and after the Russian invasion of Ukraine began in February 2022, collaboration between the JINR and other institutes completely ceased due to sanctions. Consequently, the JINR now plans to try the 248Cm+54Cr reaction instead. A preparatory experiment for the use of 54Cr projectiles was conducted in late 2023, successfully synthesising 288Lv in the 238U+54Cr reaction, and the hope is for experiments to synthesise element 120 to begin by 2025. Starting from 2022, plans have also been made to use 88-inch cyclotron in the Lawrence Berkeley National Laboratory (LBNL) in Berkeley, California, United States to attempt to make new elements using 50Ti projectiles. First, the 244Pu+50Ti reaction was tested, successfully creating two atoms of 290Lv in 2024. Since this was successful, an attempt to make element 120 in the 249Cf+50Ti reaction is planned to begin in 2025. The Lawrence Livermore National Laboratory (LLNL), which previously collaborated with the JINR, will collaborate with the LBNL on this project. Unbiunium (E121) The synthesis of element 121 (unbiunium) was first attempted in 1977 by bombarding a target of uranium-238 with copper-65 ions at the Gesellschaft für Schwerionenforschung in Darmstadt, Germany: + → 303121* → no atoms No atoms were identified. Unbibium (E122) The first attempts to synthesize element 122 (unbibium) were performed in 1972 by Flerov et al. at the Joint Institute for Nuclear Research (JINR), using the heavy-ion induced hot fusion reactions: + → 304, 306122* → no atoms These experiments were motivated by early predictions on the existence of an island of stability at N = 184 and Z > 120. No atoms were detected and a yield limit of 5 nb (5,000 pb) was measured. Current results (see flerovium) have shown that the sensitivity of these experiments were too low by at least 3 orders of magnitude. In 2000, the Gesellschaft für Schwerionenforschung (GSI) Helmholtz Center for Heavy Ion Research performed a very similar experiment with much higher sensitivity: + → 308122* → no atoms These results indicate that the synthesis of such heavier elements remains a significant challenge and further improvements of beam intensity and experimental efficiency is required. The sensitivity should be increased to 1 fb in the future for better quality results. Another unsuccessful attempt to synthesize element 122 was carried out in 1978 at the GSI Helmholtz Center, where a natural erbium target was bombarded with xenon-136 ions: + → 298, 300, 302, 303, 304, 306122* → no atoms In particular, the reaction between 170Er and 136Xe was expected to yield alpha-emitters with half-lives of microseconds that would decay down to isotopes of flerovium with half-lives perhaps increasing up to several hours, as flerovium is predicted to lie near the center of the island of stability. After twelve hours of irradiation, nothing was found in this reaction. Following a similar unsuccessful attempt to synthesize element 121 from 238U and 65Cu, it was concluded that half-lives of superheavy nuclei must be less than one microsecond or the cross sections are very small. More recent research into synthesis of superheavy elements suggests that both conclusions are true. The two attempts in the 1970s to synthesize element 122 were both propelled by the research investigating whether superheavy elements could potentially be naturally occurring. Several experiments studying the fission characteristics of various superheavy compound nuclei such as 306122* were performed between 2000 and 2004 at the Flerov Laboratory of Nuclear Reactions. Two nuclear reactions were used, namely 248Cm + 58Fe and 242Pu + 64Ni. The results reveal how superheavy nuclei fission predominantly by expelling closed shell nuclei such as 132Sn (Z = 50, N = 82). It was also found that the yield for the fusion-fission pathway was similar between 48Ca and 58Fe projectiles, suggesting a possible future use of 58Fe projectiles in superheavy element formation. Unbiquadium (E124) Scientists at GANIL (Grand Accélérateur National d'Ions Lourds) attempted to measure the direct and delayed fission of compound nuclei of elements with Z = 114, 120, and 124 in order to probe shell effects in this region and to pinpoint the next spherical proton shell. This is because having complete nuclear shells (or, equivalently, having a magic number of protons or neutrons) would confer more stability on the nuclei of such superheavy elements, thus moving closer to the island of stability. In 2006, with full results published in 2008, the team provided results from a reaction involving the bombardment of a natural germanium target with uranium ions: + → 308, 310, 311, 312, 314124* → fission The team reported that they had been able to identify compound nuclei fissioning with half-lives > 10−18 s. This result suggests a strong stabilizing effect at Z = 124 and points to the next proton shell at Z > 120, not at Z = 114 as previously thought. A compound nucleus is a loose combination of nucleons that have not arranged themselves into nuclear shells yet. It has no internal structure and is held together only by the collision forces between the target and projectile nuclei. It is estimated that it requires around 10−14 s for the nucleons to arrange themselves into nuclear shells, at which point the compound nucleus becomes a nuclide, and this number is used by IUPAC as the minimum half-life a claimed isotope must have to potentially be recognised as being discovered. Thus, the GANIL experiments do not count as a discovery of element 124. The fission of the compound nucleus 312124 was also studied in 2006 at the tandem ALPI heavy-ion accelerator at the Laboratori Nazionali di Legnaro (Legnaro National Laboratories) in Italy: + → 312124* → fission Similarly to previous experiments conducted at the JINR (Joint Institute for Nuclear Research), fission fragments clustered around doubly magic nuclei such as 132Sn (Z = 50, N = 82), revealing a tendency for superheavy nuclei to expel such doubly magic nuclei in fission. The average number of neutrons per fission from the 312124 compound nucleus (relative to lighter systems) was also found to increase, confirming that the trend of heavier nuclei emitting more neutrons during fission continues into the superheavy mass region. Unbipentium (E125) The first and only attempt to synthesize element 125 (unbipentium) was conducted in Dubna in 19701971 using zinc ions and an americium-243 target: + → 309, 311125* → no atoms No atoms were detected, and a cross section limit of 5 nb was determined. This experiment was motivated by the possibility of greater stability for nuclei around Z ~ 126 and N ~ 184, though more recent research suggests the island of stability may instead lie at a lower atomic number (such as copernicium, Z = 112), and the synthesis of heavier elements such as element 125 will require more sensitive experiments. Unbihexium (E126) The first and only attempt to synthesize element 126 (unbihexium), which was unsuccessful, was performed in 1971 at CERN (European Organization for Nuclear Research) by René Bimbot and John M. Alexander using the hot fusion reaction: + → 316126* → no atoms High-energy (13–15 MeV) alpha particles were observed and taken as possible evidence for the synthesis of element 126. Subsequent unsuccessful experiments with higher sensitivity suggest that the 10 mb sensitivity of this experiment was too low; hence, the formation of element 126 nuclei in this reaction is highly unlikely. Unbiseptium (E127) The first and only attempt to synthesize element 127 (unbiseptium), which was unsuccessful, was performed in 1978 at the UNILAC accelerator at the GSI Helmholtz Center, where a natural tantalum target was bombarded with xenon-136 ions: + → 316, 317127* → no atoms Searches in nature A study in 1976 by a group of American researchers from several universities proposed that primordial superheavy elements, mainly livermorium, elements 124, 126, and 127, could be a cause of unexplained radiation damage (particularly radiohalos) in minerals. This prompted many researchers to search for them in nature from 1976 to 1983. A group led by Tom Cahill, a professor at the University of California at Davis, claimed in 1976 that they had detected alpha particles and X-rays with the right energies to cause the damage observed, supporting the presence of these elements. In particular, the presence of long-lived (on the order of 109 years) nuclei of elements 124 and 126, along with their decay products, at an abundance of 10−11 relative to their possible congeners uranium and plutonium, was conjectured. Others claimed that none had been detected, and questioned the proposed characteristics of primordial superheavy nuclei. In particular, they cited that any such superheavy nuclei must have a closed neutron shell at N = 184 or N = 228, and this necessary condition for enhanced stability only exists in neutron deficient isotopes of livermorium or neutron rich isotopes of the other elements that would not be beta-stable unlike most naturally occurring isotopes. This activity was also proposed to be caused by nuclear transmutations in natural cerium, raising further ambiguity upon this claimed observation of superheavy elements. On April 24, 2008, a group led by Amnon Marinov at the Hebrew University of Jerusalem claimed to have found single atoms of 292122 in naturally occurring thorium deposits at an abundance of between 10−11 and 10−12 relative to thorium. The claim of Marinov et al. was criticized by a part of the scientific community. Marinov claimed that he had submitted the article to the journals Nature and Nature Physics but both turned it down without sending it for peer review. The 292122 atoms were claimed to be superdeformed or hyperdeformed isomers, with a half-life of at least 100 million years. A criticism of the technique, previously used in purportedly identifying lighter thorium isotopes by mass spectrometry, was published in Physical Review C in 2008. A rebuttal by the Marinov group was published in Physical Review C after the published comment. A repeat of the thorium experiment using the superior method of Accelerator Mass Spectrometry (AMS) failed to confirm the results, despite a 100-fold better sensitivity. This result throws considerable doubt on the results of the Marinov collaboration with regard to their claims of long-lived isotopes of thorium, roentgenium and element 122. It is still possible that traces of unbibium might only exist in some thorium samples, although this is unlikely. The possible extent of primordial superheavy elements on Earth today is uncertain. Even if they are confirmed to have caused the radiation damage long ago, they might now have decayed to mere traces, or even be completely gone. It is also uncertain if such superheavy nuclei may be produced naturally at all, as spontaneous fission is expected to terminate the r-process responsible for heavy element formation between mass number 270 and 290, well before elements beyond 120 may be formed. A recent hypothesis tries to explain the spectrum of Przybylski's Star by naturally occurring flerovium and element 120. Predicted properties of eighth-period elements Element 118, oganesson, is the heaviest element that has been synthesized. The next two elements, elements 119 and 120, should form an 8s series and be an alkali and alkaline earth metal, respectively. Beyond element 120, the superactinide series is expected to begin, when the 8s electrons and the filling of the 8p1/2, 7d3/2, 6f, and 5g subshells determine the chemistry of these elements. Complete and accurate CCSD calculations are not available for elements beyond 122 because of the extreme complexity of the situation: the 5g, 6f, and 7d orbitals should have about the same energy level, and in the region of element 160, the 9s, 8p3/2, and 9p1/2 orbitals should also be about equal in energy. This will cause the electron shells to mix so that the block concept no longer applies very well, and will also result in novel chemical properties that will make positioning some of these elements in a periodic table very difficult. Chemical and physical properties Elements 119 and 120 {\| class="wikitable" \|+ Some predicted properties of elements 119 and 120 ! Property ! 119 ! 120 \|- ! Standard atomic weight \| [322] \| [325] \|- ! Group \| 1 \| 2 \|- ! Valence electron configuration \| 8s1 \| 8s2 \|- ! Stable oxidation states \| 1, 3 \| 2, 4 \|- ! First ionization energy \| 463.1 kJ/mol \| 563.3 kJ/mol \|- ! Metallic radius \| 260 pm \| 200 pm \|- ! Density \| 3 g/cm3 \| 7 g/cm3 \|- ! Melting point \| \| \|-sigfig= ! Boiling point \| \| \|} The first two elements of period 8 will be ununennium and unbinilium, elements 119 and 120. Their electron configurations should have the 8s orbital being filled. This orbital is relativistically stabilized and contracted; thus, elements 119 and 120 should be more like rubidium and strontium than their immediate neighbours above, francium and radium. Another effect of the relativistic contraction of the 8s orbital is that the atomic radii of these two elements should be about the same as those of francium and radium. They should behave like normal alkali and alkaline earth metals (albeit less reactive than their immediate vertical neighbours), normally forming +1 and +2 oxidation states, respectively, but the relativistic destabilization of the 7p3/2 subshell and the relatively low ionization energies of the 7p3/2 electrons should make higher oxidation states like +3 and +4 (respectively) possible as well. Superactinides The superactinides may range from elements 121 through 157, which can be classified as the 5g and 6f elements of the eighth period, together with the first 7d element. In the superactinide series, the 7d, 8p, 6f and 5g shells should all fill simultaneously. This creates very complicated situations, so much so that complete and accurate CCSD calculations have been done only for elements 121 and 122. The first superactinide, unbiunium (element 121), should be similar to lanthanum and actinium: its main oxidation state should be +3, although the closeness of the valence subshells' energy levels may permit higher oxidation states, just as in elements 119 and 120. Relativistic stabilization of the 8p subshell should result in a ground-state 8s8p valence electron configuration for element 121, in contrast to the ds configurations of lanthanum and actinium; nevertheless, this anomalous configuration does not appear to affect its calculated chemistry, which remains similar to that of actinium. Its first ionization energy is predicted to be 429.4 kJ/mol, which would be lower than those of all known elements except for the alkali metals potassium, rubidium, caesium, and francium: this value is even lower than that of the period 8 alkali metal ununennium (463.1 kJ/mol). Similarly, the next superactinide, unbibium (element 122), may be similar to cerium and thorium, with a main oxidation state of +4, but would have a ground-state 7d8s8p or 8s8p valence electron configuration, unlike thorium's 6d7s configuration. Hence, its first ionization energy would be smaller than thorium's (Th: 6.3 eV; element 122: 5.6 eV) because of the greater ease of ionizing unbibium's 8p electron than thorium's 6d electron. The collapse of the 5g orbital itself is delayed until around element 125; the electron configurations of the 119-electron isoelectronic series are expected to be [Og]8s for elements 119 through 122, [Og]6f for elements 123 and 124, and [Og]5g for element 125 onwards. In the first few superactinides, the binding energies of the added electrons are predicted to be small enough that they can lose all their valence electrons; for example, unbihexium (element 126) could easily form a +8 oxidation state, and even higher oxidation states for the next few elements may be possible. Element 126 is also predicted to display a variety of other oxidation states: recent calculations have suggested a stable monofluoride 126F may be possible, resulting from a bonding interaction between the 5g orbital on element 126 and the 2p orbital on fluorine. Other predicted oxidation states include +2, +4, and +6; +4 is expected to be the most usual oxidation state of unbihexium. The superactinides from unbipentium (element 125) to unbiennium (element 129) are predicted to exhibit a +6 oxidation state and form hexafluorides, though 125F and 126F are predicted to be relatively weakly bound. The bond dissociation energies are expected to greatly increase at element 127 and even more so at element 129. This suggests a shift from strong ionic character in fluorides of element 125 to more covalent character, involving the 8p orbital, in fluorides of element 129. The bonding in these superactinide hexafluorides is mostly between the highest 8p subshell of the superactinide and the 2p subshell of fluorine, unlike how uranium uses its 5f and 6d orbitals for bonding in uranium hexafluoride. Despite the ability of early superactinides to reach high oxidation states, it has been calculated that the 5g electrons will be most difficult to ionize; the 125 and 126 ions are expected to bear a 5g configuration, similar to the 5f configuration of the Np ion. Similar behavior is observed in the low chemical activity of the 4f electrons in lanthanides; this is a consequence of the 5g orbitals being small and deeply buried in the electron cloud. The presence of electrons in g-orbitals, which do not exist in the ground state electron configuration of any currently known element, should allow presently unknown hybrid orbitals to form and influence the chemistry of the superactinides in new ways, although the absence of g electrons in known elements makes predicting superactinide chemistry more difficult. {\| class="wikitable" \|+ Some predicted compounds of the superactinides (X = a halogen) ! ! 121 ! 122 ! 123 ! 124 ! 125 ! 126 ! 127 ! 128 ! 129 ! 132 ! 142 ! 143 ! 144 ! 145 ! 146 ! 148 ! 153 ! 154 ! 155 ! 156 ! 157 \|- ! Compound \| 121X3 \| 122X4 \| 123X5 \| 124X6 \| 125F125F6 \| 126F126F6126O4 \| 127F6 \| 128F6 \| 129F129F6 \| \| 142X4142X6 \| 143F6 \| 144X6144F8144O4 \| 145F6 \| \| 148O6 \| \| \| \| \| \|- ! Analogs \| LaX3AcX3 \| CeX4ThX4 \| \| \| \| \| \| \| \| \| ThF4 \| \| UF6PuF8PuO4 \| \| \| UO6 \| \| \| \| \| \|- ! Oxidation states \| 3 \| 4 \| 5 \| 6 \| 1, 6, 7 \| 1, 2, 4, 6, 8 \| 6 \| 6 \| 1, 6 \| 6 \| 4, 6 \| 6, 8 \| 3, 4, 5, 6, 8 \| 6 \| 8 \| 12 \| 3 \| 0, 2 \| 3, 5 \| 2 \| 3 \|} In the later superactinides, the oxidation states should become lower. By element 132, the predominant most stable oxidation state will be only +6; this is further reduced to +3 and +4 by element 144, and at the end of the superactinide series it will be only +2 (and possibly even 0) because the 6f shell, which is being filled at that point, is deep inside the electron cloud and the 8s and 8p electrons are bound too strongly to be chemically active. The 5g shell should be filled at element 144 and the 6f shell at around element 154, and at this region of the superactinides the 8p electrons are bound so strongly that they are no longer active chemically, so that only a few electrons can participate in chemical reactions. Calculations by Fricke et al. predict that at element 154, the 6f shell is full and there are no d- or other electron wave functions outside the chemically inactive 8s and 8p1/2 shells. This may cause element 154 to be rather unreactive with noble gas-like properties. Calculations by Pyykkö nonetheless expect that at element 155, the 6f shell is still chemically ionizable: 155 should have a full 6f shell, and the fourth ionization potential should be between those of terbium and dysprosium, both of which are known in the +4 state. Similarly to the lanthanide and actinide contractions, there should be a superactinide contraction in the superactinide series where the ionic radii of the superactinides are smaller than expected. In the lanthanides, the contraction is about 4.4 pm per element; in the actinides, it is about 3 pm per element. The contraction is larger in the lanthanides than in the actinides due to the greater localization of the 4f wave function as compared to the 5f wave function. Comparisons with the wave functions of the outer electrons of the lanthanides, actinides, and superactinides lead to a prediction of a contraction of about 2 pm per element in the superactinides; although this is smaller than the contractions in the lanthanides and actinides, its total effect is larger due to the fact that 32 electrons are filled in the deeply buried 5g and 6f shells, instead of just 14 electrons being filled in the 4f and 5f shells in the lanthanides and actinides, respectively. Pekka Pyykkö divides these superactinides into three series: a 5g series (elements 121 to 138), an 8p1/2 series (elements 139 to 140), and a 6f series (elements 141 to 155), also noting that there would be a great deal of overlapping between energy levels and that the 6f, 7d, or 8p1/2 orbitals could well also be occupied in the early superactinide atoms or ions. He also expects that they would behave more like "superlanthanides", in the sense that the 5g electrons would mostly be chemically inactive, similarly to how only one or two 4f electrons in each lanthanide are ever ionized in chemical compounds. He also predicted that the possible oxidation states of the superactinides might rise very high in the 6f series, to values such as +12 in element 148. Andrey Kulsha has called the elements 121 to 156 "ultransition" elements and has proposed to split them into two series of eighteen each, one from elements 121 to 138 and another from elements 139 to 156. The first would be analogous to the lanthanides, with oxidation states mainly ranging from +4 to +6, as the filling of the 5g shell dominates and neighbouring elements are very similar to each other, creating an analogy to uranium, neptunium, and plutonium. The second would be analogous to the actinides: at the beginning (around elements in the 140s) very high oxidation states would be expected as the 6f shell rises above the 7d one, but after that the typical oxidation states would lower and in elements in the 150s onwards the 8p electrons would stop being chemically active. Because the two rows are separated by the addition of a complete 5g subshell, they could be considered analogues of each other as well. As an example from the late superactinides, element 156 is expected to exhibit mainly the +2 oxidation state, on account of its electron configuration with easily removed 7d electrons over a stable [Og]5g6f8s8p core. It can thus be considered a heavier congener of nobelium, which likewise has a pair of easily removed 7s electrons over a stable [Rn]5f core, and is usually in the +2 state (strong oxidisers are required to obtain nobelium in the +3 state). Its first ionization energy should be about 400 kJ/mol and its metallic radius approximately 170 picometers. With a relative atomic mass of around 445 u, it should be a very heavy metal with a density of around 26 g/cm3. Elements 157 to 166 The 7d transition metals in period 8 are expected to be elements 157 to 166. Although the 8s and 8p1/2 electrons are bound so strongly in these elements that they should not be able to take part in any chemical reactions, the 9s and 9p1/2 levels are expected to be readily available for hybridization. These 7d elements should be similar to the 4d elements yttrium through cadmium. In particular, element 164 with a 7d109s0 electron configuration shows clear analogies with palladium with its 4d105s0 electron configuration. The noble metals of this series of transition metals are not expected to be as noble as their lighter homologues, due to the absence of an outer s shell for shielding and also because the 7d shell is strongly split into two subshells due to relativistic effects. This causes the first ionization energies of the 7d transition metals to be smaller than those of their lighter congeners. Theoretical interest in the chemistry of unhexquadium is largely motivated by theoretical predictions that it, especially the isotopes 472164 and 482164 (with 164 protons and 308 or 318 neutrons), would be at the center of a hypothetical second island of stability (the first being centered on copernicium, particularly the isotopes 291Cn, 293Cn, and 296Cn which are expected to have half-lives of centuries or millennia). Calculations predict that the 7d electrons of element 164 (unhexquadium) should participate very readily in chemical reactions, so that it should be able to show stable +6 and +4 oxidation states in addition to the normal +2 state in aqueous solutions with strong ligands. Element 164 should thus be able to form compounds like 164(CO)4, 164(PF3)4 (both tetrahedral like the corresponding palladium compounds), and (linear), which is very different behavior from that of lead, which element 164 would be a heavier homologue of if not for relativistic effects. Nevertheless, the divalent state would be the main one in aqueous solution (although the +4 and +6 states would be possible with stronger ligands), and unhexquadium(II) should behave more similarly to lead than unhexquadium(IV) and unhexquadium(VI). Element 164 is expected to be a soft Lewis acid and have Ahrlands softness parameter close to 4 eV. It should be at most moderately reactive, having a first ionization energy that should be around 685 kJ/mol, comparable to that of molybdenum. Due to the lanthanide, actinide, and superactinide contractions, element 164 should have a metallic radius of only 158 pm, very close to that of the much lighter magnesium, despite its expected atomic weight of around 474 u which is about 19.5 times the atomic weight of magnesium. This small radius and high weight cause it to be expected to have an extremely high density of around 46 g·cm−3, over twice that of osmium, currently the most dense element known, at 22.61 g·cm−3; element 164 should be the second most dense element in the first 172 elements in the periodic table, with only its neighbor unhextrium (element 163) being more dense (at 47 g·cm−3). Metallic element 164 should have a very large cohesive energy (enthalpy of crystallization) due to its covalent bonds, most probably resulting in a high melting point. In the metallic state, element 164 should be quite noble and analogous to palladium and platinum. Fricke et al. suggested some formal similarities to oganesson, as both elements have closed-shell configurations and similar ionisation energies, although they note that while oganesson would be a very bad noble gas, element 164 would be a good noble metal. Elements 165 (unhexpentium) and 166 (unhexhexium), the last two 7d metals, should behave similarly to alkali and alkaline earth metals when in the +1 and +2 oxidation states, respectively. The 9s electrons should have ionization energies comparable to those of the 3s electrons of sodium and magnesium, due to relativistic effects causing the 9s electrons to be much more strongly bound than non-relativistic calculations would predict. Elements 165 and 166 should normally exhibit the +1 and +2 oxidation states, respectively, although the ionization energies of the 7d electrons are low enough to allow higher oxidation states like +3 for element 165. The oxidation state +4 for element 166 is less likely, creating a situation similar to the lighter elements in groups 11 and 12 (particularly gold and mercury). As with mercury but not copernicium, ionization of element 166 to 1662+ is expected to result in a 7d10 configuration corresponding to the loss of the s-electrons but not the d-electrons, making it more analogous to the lighter "less relativistic" group 12 elements zinc, cadmium, and mercury. {\| class="wikitable" \|+ Some predicted properties of elements 156–166The metallic radii and densities are first approximations.Most analogous group is given first, followed by other similar groups. ! Property ! 156 ! 157 ! 158 ! 159 ! 160 ! 161 ! 162 ! 163 ! 164 ! 165 ! 166 \|- ! Standard atomic weight \| [445] \| [448] \| [452] \| [456] \| [459] \| [463] \| [466] \| [470] \| [474] \| [477] \| [481] \|- ! Group \| Yb group \| 3 \| 4 \| 5 \| 6 \| 7 \| 8 \| 9 \| 10 \| 11(1) \| 12(2) \|- ! Valence electron configuration \| 7d2 \| 7d3 \| 7d4 \| 7d5 \| 7d6 \| 7d7 \| 7d8 \| 7d9 \| 7d10 \| 7d10 9s1 \| 7d10 9s2 \|- ! Stable oxidation states \| 2 \| 3 \| 4 \| 1, 5 \| 2, 6 \| 3, 7 \| 4, 8 \| 5 \| 0, 2, 4, 6 \| 1, 3 \| 2 \|- ! First ionization energy \| 400 kJ/mol \| 450 kJ/mol \| 520 kJ/mol \| 340 kJ/mol \| 420 kJ/mol \| 470 kJ/mol \| 560 kJ/mol \| 620 kJ/mol \| 690 kJ/mol \| 520 kJ/mol \| 630 kJ/mol \|- ! Metallic radius \| 170 pm \| 163 pm \| 157 pm \| 152 pm \| 148 pm \| 148 pm \| 149 pm \| 152 pm \| 158 pm \| 250 pm \| 200 pm \|- ! Density \| 26 g/cm3 \| 28 g/cm3 \| 30 g/cm3 \| 33 g/cm3 \| 36 g/cm3 \| 40 g/cm3 \| 45 g/cm3 \| 47 g/cm3 \| 46 g/cm3 \| 7 g/cm3 \| 11 g/cm3 \|} Elements 167 to 172 The next six elements on the periodic table are expected to be the last main-group elements in their period, and are likely to be similar to the 5p elements indium through xenon. In elements 167 to 172, the 9p1/2 and 8p3/2 shells will be filled. Their energy eigenvalues are so close together that they behave as one combined p-subshell, similar to the non-relativistic 2p and 3p subshells. Thus, the inert-pair effect does not occur and the most common oxidation states of elements 167 to 170 are expected to be +3, +4, +5, and +6, respectively. Element 171 (unseptunium) is expected to show some similarities to the halogens, showing various oxidation states ranging from −1 to +7, although its physical properties are expected to be closer to that of a metal. Its electron affinity is expected to be 3.0 eV, allowing it to form H171, analogous to a hydrogen halide. The 171− ion is expected to be a soft base, comparable to iodide (I−). Element 172 (unseptbium) is expected to be a noble gas with chemical behaviour similar to that of xenon, as their ionization energies should be very similar (Xe, 1170.4 kJ/mol; element 172, 1090 kJ/mol). The only main difference between them is that element 172, unlike xenon, is expected to be a liquid or a solid at standard temperature and pressure due to its much higher atomic weight. Unseptbium is expected to be a strong Lewis acid, forming fluorides and oxides, similarly to its lighter congener xenon. Because of some analogy of elements 165–172 to periods 2 and 3, Fricke et al. considered them to form a ninth period of the periodic table, while the eighth period was taken by them to end at the noble metal element 164. This ninth period would be similar to the second and third period in having no transition metals. That being said, the analogy is incomplete for elements 165 and 166; although they do start a new s-shell (9s), this is above a d-shell, making them chemically more similar to groups 11 and 12. {\| class="wikitable" \|+ Some predicted properties of elements 167–172The metallic or covalent radii and densities are first approximations. ! Property ! 167 ! 168 ! 169 ! 170 ! 171 ! 172 \|- ! Standard atomic weight \| [485] \| [489] \| [493] \| [496] \| [500] \| [504] \|- ! Group \| 13 \| 14 \| 15 \| 16 \| 17 \| 18 \|- ! Valence electron configuration \| 9s2 9p1 \| 9s2 9p2 \| 9s2 9p2 8p1 \| 9s2 9p2 8p2 \| 9s2 9p2 8p3 \| 9s2 9p2 8p4 \|- ! Stable oxidation states \| 3 \| 4 \| 5 \| 6 \| −1, 3, 7 \| 0, 4, 6, 8 \|- ! First ionization energy \| 620 kJ/mol \| 720 kJ/mol \| 800 kJ/mol \| 890 kJ/mol \| 984 kJ/mol \| 1090 kJ/mol \|- ! Metallic or covalent radius \| 190 pm \| 180 pm \| 175 pm \| 170 pm \| 165 pm \| 220 pm \|- ! Density \| 17 g/cm3 \| 19 g/cm3 \| 18 g/cm3 \| 17 g/cm3 \| 16 g/cm3 \| 9 g/cm3 \|} Beyond element 172 Beyond element 172, there is the potential to fill the 6g, 7f, 8d, 10s, 10p1/2, and perhaps 6h11/2 shells. These electrons would be very loosely bound, potentially rendering extremely high oxidation states reachable, though the electrons would become more tightly bound as the ionic charge rises. Thus, there will probably be another very long transition series, like the superactinides. In element 173 (unsepttrium), the outermost electron might enter the 6g7/2, 9p3/2, or 10s subshells. Because spin–orbit interactions would create a very large energy gap between these and the 8p3/2 subshell, this outermost electron is expected to be very loosely bound and very easily lost to form a 173+ cation. As a result, element 173 is expected to behave chemically like an alkali metal, and one that might be far more reactive than even caesium (francium and element 119 being less reactive than caesium due to relativistic effects): the calculated ionisation energy for element 173 is 3.070 eV, compared to the experimentally known 3.894 eV for caesium. Element 174 (unseptquadium) may add an 8d electron and form a closed-shell 1742+ cation; its calculated ionisation energy is 3.614 eV. Element 184 (unoctquadium) was significantly targeted in early predictions, as it was originally speculated that 184 would be a proton magic number: it is predicted to have an electron configuration of [172] 6g5 7f4 8d3, with at least the 7f and 8d electrons chemically active. Its chemical behaviour is expected to be similar to uranium and neptunium, as further ionisation past the +6 state (corresponding to removal of the 6g electrons) is likely to be unprofitable; the +4 state should be most common in aqueous solution, with +5 and +6 reachable in solid compounds. End of the periodic table The number of physically possible elements is unknown. A low estimate is that the periodic table may end soon after the island of stability, which is expected to center on Z = 126, as the extension of the periodic and nuclide tables is restricted by the proton and the neutron drip lines and stability toward alpha decay and spontaneous fission. One calculation by Y. Gambhir et al., analyzing nuclear binding energy and stability in various decay channels, suggests a limit to the existence of bound nuclei at Z = 146. Other predictions of an end to the periodic table include Z = 128 (John Emsley) and Z = 155 (Albert Khazan). Elements above the atomic number 137 It is a "folk legend" among physicists that Richard Feynman suggested that neutral atoms could not exist for atomic numbers greater than Z = 137, on the grounds that the relativistic Dirac equation predicts that the ground-state energy of the innermost electron in such an atom would be an imaginary number. Here, the number 137 arises as the inverse of the fine-structure constant. By this argument, neutral atoms cannot exist beyond atomic number 137, and therefore a periodic table of elements based on electron orbitals breaks down at this point. However, this argument presumes that the atomic nucleus is pointlike. A more accurate calculation must take into account the small, but nonzero, size of the nucleus, which is predicted to push the limit further to Z ≈ 173. Bohr model The Bohr model exhibits difficulty for atoms with atomic number greater than 137, for the speed of an electron in a 1s electron orbital, v, is given by where Z is the atomic number, and α is the fine-structure constant, a measure of the strength of electromagnetic interactions. Under this approximation, any element with an atomic number of greater than 137 would require 1s electrons to be traveling faster than c, the speed of light. Hence, the non-relativistic Bohr model is inaccurate when applied to such an element. Relativistic Dirac equation The relativistic Dirac equation gives the ground state energy as where m is the rest mass of the electron. For Z > 137, the wave function of the Dirac ground state is oscillatory, rather than bound, and there is no gap between the positive and negative energy spectra, as in the Klein paradox. More accurate calculations taking into account the effects of the finite size of the nucleus indicate that the binding energy first exceeds 2mc2 for Z > Zcr probably between 168 and 172. For Z > Zcr, if the innermost orbital (1s) is not filled, the electric field of the nucleus will pull an electron out of the vacuum, resulting in the spontaneous emission of a positron. This diving of the 1s subshell into the negative continuum has often been taken to constitute an "end" to the periodic table, but in fact it does not impose such a limit, as such resonances can be interpreted as Gamow states. Nonetheless, the accurate description of such states in a multi-electron system, needed to extend calculations and the periodic table past Zcr ≈ 172, are still open problems. Atoms with atomic numbers above Zcr ≈ 172 have been termed supercritical atoms. Supercritical atoms cannot be totally ionised because their 1s subshell would be filled by spontaneous pair creation in which an electron-positron pair is created from the negative continuum, with the electron being bound and the positron escaping. However, the strong field around the atomic nucleus is restricted to a very small region of space, so that the Pauli exclusion principle forbids further spontaneous pair creation once the subshells that have dived into the negative continuum are filled. Elements 173–184 have been termed weakly supercritical atoms as for them only the 1s shell has dived into the negative continuum; the 2p1/2 shell is expected to join around element 185 and the 2s shell around element 245. Experiments have so far not succeeded in detecting spontaneous pair creation from assembling supercritical charges through the collision of heavy nuclei (e.g. colliding lead with uranium to momentarily give an effective Z of 174; uranium with uranium gives effective Z = 184 and uranium with californium gives effective Z = 190). Even though passing Zcr does not mean elements can no longer exist, the increasing concentration of the 1s density close to the nucleus would likely make these electrons more vulnerable to K electron capture as Zcr is approached. For such heavy elements, these 1s electrons would likely spend a significant fraction of time so close to the nucleus that they are actually inside it. This may pose another limit to the periodic table. Because of the factor of m, muonic atoms become supercritical at a much larger atomic number of around 2200, as muons are about 207 times as heavy as electrons. Quark matter It has also been posited that in the region beyond A > 300, an entire "continent of stability" consisting of a hypothetical phase of stable quark matter, comprising freely flowing up and down quarks rather than quarks bound into protons and neutrons, may exist. Such a form of matter is theorized to be a ground state of baryonic matter with a greater binding energy per baryon than nuclear matter, favoring the decay of nuclear matter beyond this mass threshold into quark matter. If this state of matter exists, it could possibly be synthesized in the same fusion reactions leading to normal superheavy nuclei, and would be stabilized against fission as a consequence of its stronger binding that is enough to overcome Coulomb repulsion. Calculations published in 2020 suggest stability of up-down quark matter (udQM) nuggets against conventional nuclei beyond A ~ 266, and also show that udQM nuggets become supercritical earlier (Zcr ~ 163, A ~ 609) than conventional nuclei (Zcr ~ 177, A ~ 480). Nuclear properties Magic numbers and the island of stability The stability of nuclei decreases greatly with the increase in atomic number after curium, element 96, so that all isotopes with an atomic number above 101 decay radioactively with a half-life under a day. No elements with atomic numbers above 82 (after lead) have stable isotopes. Nevertheless, because of reasons not very well understood yet, there is a slight increased nuclear stability around atomic numbers 110–114, which leads to the appearance of what is known in nuclear physics as the "island of stability". This concept, proposed by University of California professor Glenn Seaborg, explains why superheavy elements last longer than predicted. Calculations according to the Hartree–Fock–Bogoliubov method using the non-relativistic Skyrme interaction have proposed Z = 126 as a closed proton shell. In this region of the periodic table, N = 184, N = 196, and N = 228 have been suggested as closed neutron shells. Therefore, the isotopes of most interest are 310126, 322126, and 354126, for these might be considerably longer-lived than other isotopes. Element 126, having a magic number of protons, is predicted to be more stable than other elements in this region, and may have nuclear isomers with very long half-lives. It is also possible that the island of stability is instead centered at 306122, which may be spherical and doubly magic. Probably, the island of stability occurs around Z = 114–126 and N = 184, with lifetimes probably around hours to days. Beyond the shell closure at N = 184, spontaneous fission lifetimes should drastically drop below 10−15 seconds – too short for a nucleus to obtain an electron cloud and participate in any chemistry. That being said, such lifetimes are very model-dependent, and predictions range across many orders of magnitude. Taking nuclear deformation and relativistic effects into account, an analysis of single-particle levels predicts new magic numbers for superheavy nuclei at Z = 126, 138, 154, and 164 and N = 228, 308, and 318. Therefore, in addition to the island of stability centered at 291Cn, 293Cn, and 298Fl, further islands of stability may exist around the doubly magic 354126 as well as 472164 or 482164. These nuclei are predicted to be beta-stable and decay by alpha emission or spontaneous fission with relatively long half-lives, and confer additional stability on neighboring N = 228 isotones and elements 152–168, respectively. On the other hand, the same analysis suggests that proton shell closures may be relatively weak or even nonexistent in some cases such as 354126, meaning that such nuclei might not be doubly magic and stability will instead be primarily determined by strong neutron shell closures. Additionally, due to the enormously greater forces of electromagnetic repulsion that must be overcome by the strong force at the second island (Z = 164), it is possible that nuclei around this region only exist as resonances and cannot stay together for a meaningful amount of time. It is also possible that some of the superactinides between these series may not actually exist because they are too far from both islands, in which case the periodic table might end around Z = 130. The area of elements 121–156 where periodicity is in abeyance is quite similar to the gap between the two islands. Beyond element 164, the fissility line defining the limit of stability with respect to spontaneous fission may converge with the neutron drip line, posing a limit to the existence of heavier elements. Nevertheless, further magic numbers have been predicted at Z = 210, 274, and 354 and N = 308, 406, 524, 644, and 772, with two beta-stable doubly magic nuclei found at 616210 and 798274; the same calculation method reproduced the predictions for 298Fl and 472164. (The doubly magic nuclei predicted for Z = 354 are beta-unstable, with 998354 being neutron-deficient and 1126354 being neutron-rich.) Although additional stability toward alpha decay and fission are predicted for 616210 and 798274, with half-lives up to hundreds of microseconds for 616210, there will not exist islands of stability as significant as those predicted at Z = 114 and 164. As the existence of superheavy elements is very strongly dependent on stabilizing effects from closed shells, nuclear instability and fission will likely determine the end of the periodic table beyond these islands of stability. The International Union of Pure and Applied Chemistry (IUPAC) defines an element to exist if its lifetime is longer than 10−14 seconds, which is the time it takes for the nucleus to form an electron cloud. However, a nuclide is generally considered to exist if its lifetime is longer than about 10−22 seconds, which is the time it takes for nuclear structure to form. Consequently, it is possible that some Z values can only be realised in nuclides and that the corresponding elements do not exist. It is also possible that no further islands actually exist beyond 126, as the nuclear shell structure gets smeared out (as the electron shell structure already is expected to be around oganesson) and low-energy decay modes become readily available. In some regions of the table of nuclides, there are expected to be additional regions of stability due to non-spherical nuclei that have different magic numbers than spherical nuclei do; the egg-shaped 270Hs is one such deformed doubly magic nucleus. In the superheavy region, the strong Coulomb repulsion of protons may cause some nuclei, including isotopes of oganesson, to assume a bubble shape in the ground state with a reduced central density of protons, unlike the roughly uniform distribution inside most smaller nuclei. Such a shape would have a very low fission barrier, however. Even heavier nuclei in some regions, such as 342136 and 466156, may instead become toroidal or red blood cell-like in shape, with their own magic numbers and islands of stability, but they would also fragment easily. Predicted decay properties of undiscovered elements As the main island of stability is thought to lie around 291Cn and 293Cn, undiscovered elements beyond oganesson may be very unstable and undergo alpha decay or spontaneous fission in microseconds or less. The exact region in which half-lives exceed one microsecond is unknown, though various models suggest that isotopes of elements heavier than unbinilium that may be produced in fusion reactions with available targets and projectiles will have half-lives under one microsecond and therefore may not be detected. It is consistently predicted that there will exist regions of stability at N = 184 and N = 228, and possibly also at Z ~ 124 and N ~ 198. These nuclei may have half-lives of a few seconds and undergo predominantly alpha decay and spontaneous fission, though minor beta-plus decay (or electron capture) branches may also exist. Outside these regions of enhanced stability, fission barriers are expected to drop significantly due to loss of stabilization effects, resulting in fission half-lives below 10−18 seconds, especially in even–even nuclei for which hindrance is even lower due to nucleon pairing. In general, alpha decay half-lives are expected to increase with neutron number, from nanoseconds in the most neutron-deficient isotopes to seconds closer to the beta-stability line. For nuclei with only a few neutrons more than a magic number, binding energy substantially drops, resulting in a break in the trend and shorter half-lives. The most neutron deficient isotopes of these elements may also be unbound and undergo proton emission. Cluster decay (heavy particle emission) has also been proposed as an alternative decay mode for some isotopes, posing yet another hurdle to identification of these elements. Electron configurations The following are expected electron configurations of elements 119–174 and 184. The symbol [Og] indicates the probable electron configuration of oganesson (Z = 118), which is currently the last known element. The configurations of the elements in this table are written starting with [Og] because oganesson is expected to be the last prior element with a closed-shell (inert gas) configuration, 1s2 2s2 2p6 3s2 3p6 3d10 4s2 4p6 4d10 4f14 5s2 5p6 5d10 5f14 6s2 6p6 6d10 7s2 7p6. Similarly, the [172] in the configurations for elements 173, 174, and 184 denotes the likely closed-shell configuration of element 172. Beyond element 123, no complete calculations are available and hence the data in this table must be taken as tentative. In the case of element 123, and perhaps also heavier elements, several possible electron configurations are predicted to have very similar energy levels, such that it is very difficult to predict the ground state. All configurations that have been proposed (since it was understood that the Madelung rule probably stops working here) are included. The predicted block assignments up to 172 are Kulsha's, following the expected available valence orbitals. There is, however, not a consensus in the literature as to how the blocks should work after element 138. {\| class="wikitable" ! colspan="3" \| Chemical element !! Block !! Predicted electron configurations \|-bgcolor="" \|\| 119 \|\| Uue \|\| Ununennium \|\|s-block \|\|[Og] 8s1 \|-bgcolor="" \|\| 120 \|\| Ubn \|\| Unbinilium \|\|s-block \|\|[Og] 8s2 \|-bgcolor="" \|\| 121 \|\| Ubu \|\| Unbiunium \|\|g-block \|\| [Og] 8s2 8p \|-bgcolor="" \|\| 122 \|\| Ubb \|\| Unbibium \|\|g-block \|\| [Og] 8s2 8p[Og] 7d1 8s2 8p \|-bgcolor="" \|\| 123 \|\| Ubt \|\| Unbitrium \|\|g-block \|\| [Og] 6f1 8s2 8p[Og] 6f1 7d1 8s2 8p[Og] 6f2 8s2 8p[Og] 8s2 8p 8p \|-bgcolor="" \|\| 124 \|\| Ubq \|\| Unbiquadium \|\|g-block \|\| [Og] 6f2 8s2 8p[Og] 6f3 8s2 8p \|-bgcolor="" \|\| 125 \|\| Ubp \|\| Unbipentium \|\|g-block \|\| [Og] 6f4 8s2 8p[Og] 5g1 6f2 8s2 8p[Og] 5g1 6f3 8s2 8p[Og] 8s2 0.81(5g1 6f2 8p) + 0.17(5g1 6f1 7d2 8p) + 0.02(6f3 7d1 8p) \|-bgcolor="" \|\| 126 \|\| Ubh \|\| Unbihexium \|\|g-block \|\| [Og] 5g1 6f4 8s2 8p[Og] 5g2 6f2 8s2 8p[Og] 5g2 6f3 8s2 8p[Og] 8s2 0.998(5g2 6f3 8p) + 0.002(5g2 6f2 8p) \|-bgcolor="" \|\| 127 \|\| Ubs \|\| Unbiseptium \|\|g-block \|\| [Og] 5g2 6f3 8s2 8p[Og] 5g3 6f2 8s2 8p[Og] 8s2 0.88(5g3 6f2 8p) + 0.12(5g3 6f1 7d2 8p) \|-bgcolor="" \|\| 128 \|\| Ubo \|\|Unbioctium\|\|g-block \|\| [Og] 5g3 6f3 8s2 8p[Og] 5g4 6f2 8s2 8p[Og] 8s2 0.88(5g4 6f2 8p) + 0.12(5g4 6f1 7d2 8p) \|-bgcolor="" \|\| 129 \|\| Ube \|\| Unbiennium \|\|g-block \|\| [Og] 5g4 6f3 7d1 8s2 8p[Og] 5g4 6f3 8s2 8p[Og] 5g5 6f2 8s2 8p[Og] 5g4 6f3 7d1 8s2 8p \|-bgcolor="" \|\| 130 \|\| Utn \|\| Untrinilium \|\|g-block \|\| [Og] 5g5 6f3 7d1 8s2 8p[Og] 5g5 6f3 8s2 8p[Og] 5g6 6f2 8s2 8p[Og] 5g5 6f3 7d1 8s2 8p \|-bgcolor="" \|\| 131 \|\| Utu \|\| Untriunium \|\|g-block \|\| [Og] 5g6 6f3 8s2 8p[Og] 5g7 6f2 8s2 8p[Og] 8s2 0.86(5g6 6f3 8p) + 0.14(5g6 6f2 7d2 8p) \|-bgcolor="" \|\| 132 \|\| Utb \|\| Untribium \|\|g-block \|\| [Og] 5g7 6f3 8s2 8p[Og] 5g8 6f2 8s2 8p \|-bgcolor="" \|\| 133 \|\| Utt \|\| Untritrium \|\|g-block \|\| [Og] 5g8 6f3 8s2 8p \|-bgcolor="" \|\| 134 \|\| Utq \|\| Untriquadium \|\|g-block \|\| [Og] 5g8 6f4 8s2 8p \|-bgcolor="" \|\| 135 \|\| Utp \|\| Untripentium \|\|g-block \|\| [Og] 5g9 6f4 8s2 8p \|-bgcolor="" \|\| 136 \|\| Uth \|\| Untrihexium \|\|g-block \|\| [Og] 5g10 6f4 8s2 8p \|-bgcolor="" \|\| 137 \|\| Uts \|\| Untriseptium \|\|g-block \|\| [Og] 5g11 6f4 8s2 8p \|-bgcolor="" \|\| 138 \|\| Uto \|\| Untrioctium \|\|g-block \|\| [Og] 5g12 6f4 8s2 8p[Og] 5g12 6f3 7d1 8s2 8p \|-bgcolor="" \|\| 139 \|\| Ute \|\| Untriennium \|\|g-block \|\| [Og] 5g13 6f3 7d1 8s2 8p[Og] 5g13 6f2 7d2 8s2 8p \|-bgcolor="" \|\| 140 \|\| Uqn \|\| Unquadnilium \|\|g-block \|\| [Og] 5g14 6f3 7d1 8s2 8p[Og] 5g15 6f1 8s2 8p 8p \|-bgcolor="" \|\| 141 \|\| Uqu \|\| Unquadunium \|\|g-block \|\| [Og] 5g15 6f2 7d2 8s2 8p \|-bgcolor="" \|\| 142 \|\| Uqb \|\| Unquadbium \|\|g-block \|\| [Og] 5g16 6f2 7d2 8s2 8p \|-bgcolor="" \|\| 143 \|\| Uqt \|\| Unquadtrium \|\|f-block \|\| [Og] 5g17 6f2 7d2 8s2 8p \|-bgcolor="" \|\| 144 \|\| Uqq \|\| Unquadquadium \|\|f-block \|\| [Og] 5g18 6f2 7d2 8s2 8p[Og] 5g18 6f1 7d3 8s2 8p[Og] 5g17 6f2 7d3 8s2 8p[Og] 8s2 0.95(5g17 6f2 7d3 8p) + 0.05(5g17 6f4 7d1 8p) \|-bgcolor="" \|\| 145 \|\| Uqp \|\| Unquadpentium \|\|f-block \|\| [Og] 5g18 6f3 7d2 8s2 8p \|-bgcolor="" \|\| 146 \|\| Uqh \|\| Unquadhexium \|\|f-block \|\| [Og] 5g18 6f4 7d2 8s2 8p \|-bgcolor="" \|\| 147 \|\| Uqs \|\| Unquadseptium \|\|f-block \|\| [Og] 5g18 6f5 7d2 8s2 8p \|-bgcolor="" \|\| 148 \|\| Uqo \|\| Unquadoctium \|\|f-block \|\| [Og] 5g18 6f6 7d2 8s2 8p \|-bgcolor="" \|\| 149 \|\| Uqe \|\| Unquadennium \|\|f-block \|\| [Og] 5g18 6f6 7d3 8s2 8p \|-bgcolor="" \|\| 150 \|\| Upn \|\| Unpentnilium \|\|f-block \|\| [Og] 5g18 6f6 7d4 8s2 8p[Og] 5g18 6f7 7d3 8s2 8p \|-bgcolor="" \|\| 151 \|\| Upu \|\| Unpentunium \|\|f-block \|\| [Og] 5g18 6f8 7d3 8s2 8p \|-bgcolor="" \|\| 152 \|\| Upb \|\| Unpentbium \|\|f-block \|\| [Og] 5g18 6f9 7d3 8s2 8p \|-bgcolor="" \|\| 153 \|\| Upt \|\| Unpenttrium \|\|f-block \|\| [Og] 5g18 6f10 7d3 8s2 8p[Og] 5g18 6f11 7d2 8s2 8p \|-bgcolor="" \|\| 154 \|\| Upq \|\| Unpentquadium \|\|f-block \|\| [Og] 5g18 6f11 7d3 8s2 8p[Og] 5g18 6f12 7d2 8s2 8p \|-bgcolor="" \|\| 155 \|\| Upp \|\| Unpentpentium \|\|f-block \|\| [Og] 5g18 6f12 7d3 8s2 8p[Og] 5g18 6f13 7d2 8s2 8p \|-bgcolor="" \|\| 156 \|\| Uph \|\| Unpenthexium \|\|f-block\|\| [Og] 5g18 6f13 7d3 8s2 8p[Og] 5g18 6f14 7d2 8s2 8p \|-bgcolor="" \|\| 157 \|\| Ups \|\| Unpentseptium \|\|d-block \|\| [Og] 5g18 6f14 7d3 8s2 8p \|-bgcolor="" \|\| 158 \|\| Upo \|\| Unpentoctium \|\|d-block \|\| [Og] 5g18 6f14 7d4 8s2 8p \|-bgcolor="" \|\| 159 \|\| Upe \|\| Unpentennium \|\|d-block \|\| [Og] 5g18 6f14 7d5 8s2 8p[Og] 5g18 6f14 7d4 8s2 8p 9s1 \|-bgcolor="" \|\| 160 \|\| Uhn \|\| Unhexnilium \|\|d-block \|\| [Og] 5g18 6f14 7d6 8s2 8p[Og] 5g18 6f14 7d5 8s2 8p 9s1 \|-bgcolor="" \|\| 161 \|\| Uhu \|\| Unhexunium \|\|d-block \|\| [Og] 5g18 6f14 7d7 8s2 8p[Og] 5g18 6f14 7d6 8s2 8p 9s1 \|-bgcolor="" \|\| 162 \|\| Uhb \|\| Unhexbium \|\|d-block \|\| [Og] 5g18 6f14 7d8 8s2 8p[Og] 5g18 6f14 7d7 8s2 8p 9s1 \|-bgcolor="" \|\| 163 \|\| Uht \|\| Unhextrium \|\|d-block \|\| [Og] 5g18 6f14 7d9 8s2 8p[Og] 5g18 6f14 7d8 8s2 8p 9s1 \|-bgcolor="" \|\| 164 \|\| Uhq \|\| Unhexquadium \|\|d-block \|\| [Og] 5g18 6f14 7d10 8s2 8p \|-bgcolor="" \|\| 165 \|\| Uhp \|\| Unhexpentium \|\|d-block \|\| [Og] 5g18 6f14 7d10 8s2 8p 9s1 \|-bgcolor="" \|\| 166 \|\| Uhh \|\| Unhexhexium \|\|d-block \|\|[Og] 5g18 6f14 7d10 8s2 8p 9s2 \|-bgcolor="" \|\| 167 \|\| Uhs \|\| Unhexseptium \|\|p-block \|\| [Og] 5g18 6f14 7d10 8s2 8p 9s2 9p[Og] 5g18 6f14 7d10 8s2 8p 8p 9s2 \|-bgcolor="" \|\| 168 \|\| Uho \|\| Unhexoctium \|\|p-block \|\| [Og] 5g18 6f14 7d10 8s2 8p 9s2 9p[Og] 5g18 6f14 7d10 8s2 8p 8p 9s2 \|-bgcolor="" \|\| 169 \|\| Uhe \|\| Unhexennium \|\|p-block \|\| [Og] 5g18 6f14 7d10 8s2 8p 8p 9s2 9p[Og] 5g18 6f14 7d10 8s2 8p 8p 9s2 \|-bgcolor="" \|\| 170 \|\| Usn \|\| Unseptnilium \|\|p-block \|\| [Og] 5g18 6f14 7d10 8s2 8p 8p 9s2 9p[Og] 5g18 6f14 7d10 8s2 8p 8p 9s2 \|-bgcolor="" \|\| 171 \|\| Usu \|\| Unseptunium \|\|p-block \|\| [Og] 5g18 6f14 7d10 8s2 8p 8p 9s2 9p[Og] 5g18 6f14 7d10 8s2 8p 8p 9s2 9p \|-bgcolor="" \|\| 172 \|\| \|\| Unseptbium \|\|p-block \|\| [Og] 5g18 6f14 7d10 8s2 8p 8p 9s2 9p \|- \|\| 173 \|\| Ust \|\| Unsepttrium \|\| ? \|\| [172] 6g1[172] 9p[172] 10s1 \|- \|\| 174 \|\| Usq \|\| Unseptquadium \|\| ? \|\| [172] 8d1 10s1 \|- \|\| ... \|\| ... \|\| ... \|\| ... \|\| ... \|- \|\| 184 \|\| Uoq \|\| Unoctquadium \|\| ? \|\| [172] 6g5 7f4 8d3 \|}	Physical sciences	Periods	Chemistry
68367	https://en.wikipedia.org/wiki/Computer%20chess	Computer chess	Computer chess includes both hardware (dedicated computers) and software capable of playing chess. Computer chess provides opportunities for players to practice even in the absence of human opponents, and also provides opportunities for analysis, entertainment and training. Computer chess applications that play at the level of a chess grandmaster or higher are available on hardware from supercomputers to smart phones. Standalone chess-playing machines are also available. Stockfish, Leela Chess Zero, GNU Chess, Fruit, and other free open source applications are available for various platforms. Computer chess applications, whether implemented in hardware or software, use different strategies than humans to choose their moves: they use heuristic methods to build, search and evaluate trees representing sequences of moves from the current position and attempt to execute the best such sequence during play. Such trees are typically quite large, thousands to millions of nodes. The computational speed of modern computers, capable of processing tens of thousands to hundreds of thousands of nodes or more per second, along with extension and reduction heuristics that narrow the tree to mostly relevant nodes, make such an approach effective. The first chess machines capable of playing chess or reduced chess-like games were software programs running on digital computers early in the vacuum-tube computer age (1950s). The early programs played so poorly that even a beginner could defeat them. Within 40 years, in 1997, chess engines running on super-computers or specialized hardware were capable of defeating even the best human players. By 2006, programs running on desktop PCs had attained the same capability. In 2006, Monty Newborn, Professor of Computer Science at McGill University, declared: "the science has been done". Nevertheless, solving chess is not currently possible for modern computers due to the game's extremely large number of possible variations. Computer chess was once considered the "Drosophila of AI", the edge of knowledge engineering. The field is now considered a scientifically completed paradigm, and playing chess is a mundane computing activity. Availability and playing strength In the past, stand-alone chess machines (usually microprocessors running software chess programs; occasionally specialized hardware) were sold. Today, chess engines may be installed as software on ordinary devices like smartphones and PCs, either alone or alongside GUI programs such as Chessbase and the mobile apps for Chess.com and Lichess (both primarily websites). Examples of free and open source engines include Stockfish and Leela Chess Zero (Lc0). Chess.com maintains its own proprietary engine named Torch. Some chess engines, including Stockfish, have web versions made in languages like WebAssembly and JavaScript. Most chess programs and sites offer the ability to analyze positions and games using chess engines, and some offer the ability to play against engines (which can be set to play at custom levels of strength) as though they were normal opponents. Hardware requirements for chess engines are minimal, but performance will vary with processor speed, and memory, needed to hold large transposition tables. Most modern chess engines, such as Stockfish, rely on efficiently updatable neural networks, tailored to be run exclusively on CPUs, but Lc0 uses networks reliant on GPU performance. Top engines such as Stockfish can be expected to beat the world's best players reliably, even when running on consumer-grade hardware. Types and features of chess software Perhaps the most common type of chess software are programs that simply play chess. A human player makes a move on the board, the AI calculates and plays a subsequent move, and the human and AI alternate turns until the game ends. The chess engine, which calculates the moves, and the graphical user interface (GUI) are sometimes separate programs. Different engines can be connected to the GUI, permitting play against different styles of opponent. Engines often have a simple text command-line interface, while GUIs may offer a variety of piece sets, board styles, or even 3D or animated pieces. Because recent engines are so capable, engines or GUIs may offer some way of handicapping the engine's ability, to improve the odds for a win by the human player. Universal Chess Interface (UCI) engines such Fritz or Rybka may have a built in mechanism for reducing the Elo rating of the engine (via UCI's uci_limitstrength and uci_elo parameters). Some versions of Fritz have a Handicap and Fun mode for limiting the current engine or changing the percentage of mistakes it makes or changing its style. Fritz also has a Friend Mode where during the game it tries to match the level of the player. Chess databases allow users to search through a large library of historical games, analyze them, check statistics, and formulate an opening repertoire. Chessbase (for PC) is a common program for these purposes amongst professional players, but there are alternatives such as Shane's Chess Information Database (Scid) for Windows, Mac or Linux, Chess Assistant for PC, Gerhard Kalab's Chess PGN Master for Android or Giordano Vicoli's Chess-Studio for iOS. Programs such as Playchess allow players to play against one another over the internet. Chess training programs teach chess. Chessmaster had playthrough tutorials by IM Josh Waitzkin and GM Larry Christiansen. Stefan Meyer-Kahlen offers Shredder Chess Tutor based on the Step coursebooks of Rob Brunia and Cor Van Wijgerden. Former World Champion Magnus Carlsen's Play Magnus company released a Magnus Trainer app for Android and iOS. Chessbase has Fritz and Chesster for children. Convekta provides a large number of training apps such as CT-ART and its Chess King line based on tutorials by GM Alexander Kalinin and Maxim Blokh. There is also software for handling chess problems. Computers versus humans After discovering refutation screening—the application of alpha–beta pruning to optimizing move evaluation—in 1957, a team at Carnegie Mellon University predicted that a computer would defeat the world human champion by 1967. It did not anticipate the difficulty of determining the right order to evaluate moves. Researchers worked to improve programs' ability to identify killer heuristics, unusually high-scoring moves to reexamine when evaluating other branches, but into the 1970s most top chess players believed that computers would not soon be able to play at a Master level. In 1968, International Master David Levy made a famous bet that no chess computer would be able to beat him within ten years, and in 1976 Senior Master and professor of psychology Eliot Hearst of Indiana University wrote that "the only way a current computer program could ever win a single game against a master player would be for the master, perhaps in a drunken stupor while playing 50 games simultaneously, to commit some once-in-a-year blunder". In the late 1970s chess programs suddenly began defeating highly skilled human players. The year of Hearst's statement, Northwestern University's Chess 4.5 at the Paul Masson American Chess Championship's Class B level became the first to win a human tournament. Levy won his bet in 1978 by beating Chess 4.7, but it achieved the first computer victory against a Master-class player at the tournament level by winning one of the six games. In 1980, Belle began often defeating Masters. By 1982 two programs played at Master level and three were slightly weaker. The sudden improvement without a theoretical breakthrough was unexpected, as many did not expect that Belle's ability to examine 100,000 positions a second—about eight plies—would be sufficient. The Spracklens, creators of the successful microcomputer program Sargon, estimated that 90% of the improvement came from faster evaluation speed and only 10% from improved evaluations. New Scientist stated in 1982 that computers "play terrible chess ... clumsy, inefficient, diffuse, and just plain ugly", but humans lost to them by making "horrible blunders, astonishing lapses, incomprehensible oversights, gross miscalculations, and the like" much more often than they realized; "in short, computers win primarily through their ability to find and exploit miscalculations in human initiatives". By 1982, microcomputer chess programs could evaluate up to 1,500 moves a second and were as strong as mainframe chess programs of five years earlier, able to defeat a majority of amateur players. While only able to look ahead one or two plies more than at their debut in the mid-1970s, doing so improved their play more than experts expected; seemingly minor improvements "appear to have allowed the crossing of a psychological threshold, after which a rich harvest of human error becomes accessible", New Scientist wrote. While reviewing SPOC in 1984, BYTE wrote that "Computers—mainframes, minis, and micros—tend to play ugly, inelegant chess", but noted Robert Byrne's statement that "tactically they are freer from error than the average human player". The magazine described SPOC as a "state-of-the-art chess program" for the IBM PC with a "surprisingly high" level of play, and estimated its USCF rating as 1700 (Class B). At the 1982 North American Computer Chess Championship, Monroe Newborn predicted that a chess program could become world champion within five years; tournament director and International Master Michael Valvo predicted ten years; the Spracklens predicted 15; Ken Thompson predicted more than 20; and others predicted that it would never happen. The most widely held opinion, however, stated that it would occur around the year 2000. In 1989, Levy was defeated by Deep Thought in an exhibition match. Deep Thought, however, was still considerably below World Championship level, as the reigning world champion, Garry Kasparov, demonstrated in two strong wins in 1989. It was not until a 1996 match with IBM's Deep Blue that Kasparov lost his first game to a computer at tournament time controls in Deep Blue versus Kasparov, 1996, game 1. This game was, in fact, the first time a reigning world champion had lost to a computer using regular time controls. However, Kasparov regrouped to win three and draw two of the remaining five games of the match, for a convincing victory. In May 1997, an updated version of Deep Blue defeated Kasparov 3½–2½ in a return match. A documentary mainly about the confrontation was made in 2003, titled Game Over: Kasparov and the Machine. With increasing processing power and improved evaluation functions, chess programs running on commercially available workstations began to rival top-flight players. In 1998, Rebel 10 defeated Viswanathan Anand, who at the time was ranked second in the world, by a score of 5–3. However, most of those games were not played at normal time controls. Out of the eight games, four were blitz games (five minutes plus five seconds Fischer delay for each move); these Rebel won 3–1. Two were semi-blitz games (fifteen minutes for each side) that Rebel won as well (1½–½). Finally, two games were played as regular tournament games (forty moves in two hours, one hour sudden death); here it was Anand who won ½–1½. In fast games, computers played better than humans, but at classical time controls – at which a player's rating is determined – the advantage was not so clear. In the early 2000s, commercially available programs such as Junior and Fritz were able to draw matches against former world champion Garry Kasparov and classical world champion Vladimir Kramnik. In October 2002, Vladimir Kramnik and Deep Fritz competed in the eight-game Brains in Bahrain match, which ended in a draw. Kramnik won games 2 and 3 by "conventional" anti-computer tactics – play conservatively for a long-term advantage the computer is not able to see in its game tree search. Fritz, however, won game 5 after a severe blunder by Kramnik. Game 6 was described by the tournament commentators as "spectacular". Kramnik, in a better position in the early middlegame, tried a piece sacrifice to achieve a strong tactical attack, a strategy known to be highly risky against computers who are at their strongest defending against such attacks. True to form, Fritz found a watertight defense and Kramnik's attack petered out leaving him in a bad position. Kramnik resigned the game, believing the position lost. However, post-game human and computer analysis has shown that the Fritz program was unlikely to have been able to force a win and Kramnik effectively sacrificed a drawn position. The final two games were draws. Given the circumstances, most commentators still rate Kramnik the stronger player in the match. In January 2003, Kasparov played Junior, another chess computer program, in New York City. The match ended 3–3. In November 2003, Kasparov played X3D Fritz. The match ended 2–2. In 2005, Hydra, a dedicated chess computer with custom hardware and sixty-four processors and also winner of the 14th IPCCC in 2005, defeated seventh-ranked Michael Adams 5½–½ in a six-game match (though Adams' preparation was far less thorough than Kramnik's for the 2002 series). In November–December 2006, World Champion Vladimir Kramnik played Deep Fritz. This time the computer won; the match ended 2–4. Kramnik was able to view the computer's opening book. In the first five games Kramnik steered the game into a typical "anti-computer" positional contest. He lost one game (overlooking a mate in one), and drew the next four. In the final game, in an attempt to draw the match, Kramnik played the more aggressive Sicilian Defence and was crushed. There was speculation that interest in human–computer chess competition would plummet as a result of the 2006 Kramnik-Deep Fritz match. According to Newborn, for example, "the science is done". Human–computer chess matches showed the best computer systems overtaking human chess champions in the late 1990s. For the 40 years prior to that, the trend had been that the best machines gained about 40 points per year in the Elo rating while the best humans only gained roughly 2 points per year. The highest rating obtained by a computer in human competition was Deep Thought's USCF rating of 2551 in 1988 and FIDE no longer accepts human–computer results in their rating lists. Specialized machine-only Elo pools have been created for rating machines, but such numbers, while similar in appearance, are not directly compared. In 2016, the Swedish Chess Computer Association rated computer program Komodo at 3361. Chess engines continue to improve. In 2009, chess engines running on slower hardware have reached the grandmaster level. A mobile phone won a category 6 tournament with a performance rating 2898: chess engine Hiarcs 13 running inside Pocket Fritz 4 on the mobile phone HTC Touch HD won the Copa Mercosur tournament in Buenos Aires, Argentina with 9 wins and 1 draw on August 4–14, 2009. Pocket Fritz 4 searches fewer than 20,000 positions per second. This is in contrast to supercomputers such as Deep Blue that searched 200 million positions per second. Advanced Chess is a form of chess developed in 1998 by Kasparov where a human plays against another human, and both have access to computers to enhance their strength. The resulting "advanced" player was argued by Kasparov to be stronger than a human or computer alone. This has been proven in numerous occasions, such as at Freestyle Chess events. Players today are inclined to treat chess engines as analysis tools rather than opponents. Chess grandmaster Andrew Soltis stated in 2016 "The computers are just much too good" and that world champion Magnus Carlsen won't play computer chess because "he just loses all the time and there's nothing more depressing than losing without even being in the game." Computer methods Since the era of mechanical machines that played rook and king endings and electrical machines that played other games like hex in the early years of the 20th century, scientists and theoreticians have sought to develop a procedural representation of how humans learn, remember, think and apply knowledge, and the game of chess, because of its daunting complexity, became the "Drosophila of artificial intelligence (AI)". The procedural resolution of complexity became synonymous with thinking, and early computers, even before the chess automaton era, were popularly referred to as "electronic brains". Several different schema were devised starting in the latter half of the 20th century to represent knowledge and thinking, as applied to playing the game of chess (and other games like checkers): Search based (brute force vs selective search) Search in search based schema (minimax/alpha-beta, Monte Carlo tree search) Evaluations in search based schema (machine learning, neural networks, texel tuning, genetic algorithms, gradient descent, reinforcement learning) Knowledge based (PARADISE, endgame tablebases) Using "ends-and-means" heuristics a human chess player can intuitively determine optimal outcomes and how to achieve them regardless of the number of moves necessary, but a computer must be systematic in its analysis. Most players agree that looking at least five moves ahead (ten plies) when necessary is required to play well. Normal tournament rules give each player an average of three minutes per move. On average there are more than 30 legal moves per chess position, so a computer must examine a quadrillion possibilities to look ahead ten plies (five full moves); one that could examine a million positions a second would require more than 30 years. The earliest attempts at procedural representations of playing chess predated the digital electronic age, but it was the stored program digital computer that gave scope to calculating such complexity. Claude Shannon, in 1949, laid out the principles of algorithmic solution of chess. In that paper, the game is represented by a "tree", or digital data structure of choices (branches) corresponding to moves. The nodes of the tree were positions on the board resulting from the choices of move. The impossibility of representing an entire game of chess by constructing a tree from first move to last was immediately apparent: there are an average of 36 moves per position in chess and an average game lasts about 35 moves to resignation (60-80 moves if played to checkmate, stalemate, or other draw). There are 400 positions possible after the first move by each player, about 200,000 after two moves each, and nearly 120 million after just 3 moves each. So a limited lookahead (search) to some depth, followed by using domain-specific knowledge to evaluate the resulting terminal positions was proposed. A kind of middle-ground position, given good moves by both sides, would result, and its evaluation would inform the player about the goodness or badness of the moves chosen. Searching and comparing operations on the tree were well suited to computer calculation; the representation of subtle chess knowledge in the evaluation function was not. The early chess programs suffered in both areas: searching the vast tree required computational resources far beyond those available, and what chess knowledge was useful and how it was to be encoded would take decades to discover. The developers of a chess-playing computer system must decide on a number of fundamental implementation issues. These include: Graphical user interface (GUI) – how moves are entered and communicated to the user, how the game is recorded, how the time controls are set, and other interface considerations Board representation – how a single position is represented in data structures; Search techniques – how to identify the possible moves and select the most promising ones for further examination; Leaf evaluation – how to evaluate the value of a board position, if no further search will be done from that position. Adriaan de Groot interviewed a number of chess players of varying strengths, and concluded that both masters and beginners look at around forty to fifty positions before deciding which move to play. What makes the former much better players is that they use pattern recognition skills built from experience. This enables them to examine some lines in much greater depth than others by simply not considering moves they can assume to be poor. More evidence for this being the case is the way that good human players find it much easier to recall positions from genuine chess games, breaking them down into a small number of recognizable sub-positions, rather than completely random arrangements of the same pieces. In contrast, poor players have the same level of recall for both. The equivalent of this in computer chess are evaluation functions for leaf evaluation, which correspond to the human players' pattern recognition skills, and the use of machine learning techniques in training them, such as Texel tuning, stochastic gradient descent, and reinforcement learning, which corresponds to building experience in human players. This allows modern programs to examine some lines in much greater depth than others by using forwards pruning and other selective heuristics to simply not consider moves the program assume to be poor through their evaluation function, in the same way that human players do. The only fundamental difference between a computer program and a human in this sense is that a computer program can search much deeper than a human player could, allowing it to search more nodes and bypass the horizon effect to a much greater extent than is possible with human players. Graphical user interface Computer chess programs usually support a number of common de facto standards. Nearly all of today's programs can read and write game moves as Portable Game Notation (PGN), and can read and write individual positions as Forsyth–Edwards Notation (FEN). Older chess programs often only understood long algebraic notation, but today users expect chess programs to understand standard algebraic chess notation. Starting in the late 1990s, programmers began to develop separately engines (with a command-line interface which calculates which moves are strongest in a position) or a graphical user interface (GUI) which provides the player with a chessboard they can see, and pieces that can be moved. Engines communicate their moves to the GUI using a protocol such as the Chess Engine Communication Protocol (CECP) or Universal Chess Interface (UCI). By dividing chess programs into these two pieces, developers can write only the user interface, or only the engine, without needing to write both parts of the program. (	Technology	Artificial intelligence concepts	null
68468	https://en.wikipedia.org/wiki/Dune%20buggy	Dune buggy	A dune buggy — also known as a beach buggy — is a recreational off-road vehicle with large wheels, and wide tires, designed for use on sand dunes, beaches, off-road or desert recreation. The design is usually a topless vehicle with a rear-mounted engine. A dune buggy can be created by modifying an existing vehicle or custom-building a new vehicle. Design Dune buggies are typically created by modifying an existing road vehicle, while sandrails are built from the ground up as a custom vehicle. Beetle-based buggies For dune buggies built on the chassis of a rear-engined existing vehicle, the Volkswagen Beetle has been most commonly used as the basis for the buggy, though conversions were made from other rear-engined cars (such as the Corvair and Renault Dauphine). The model is nicknamed Bug, lending partial inspiration to the term "buggy." The Beetle platform chassis was used because the rear engine layout improves traction, the air-cooled engine avoids the complexities and failure points associated with a water-cooled engine, the Beetle's front torsion bar suspension was not only considered cheap and robust, but it was also extremely easy to alter and adjust its ride-height. Furthermore, spare parts — and donor vehicles themselves — were cheap and readily available. While early dune buggy conversions were left with no body, or featured custom bodies of sheet metal (such as the EMPI Sportsters and similar buggies), glass-reinforced plastic (fiberglass) bodies, developed in the 1960s, have become the standard image of the modern buggy, and come in many shapes and sizes. The original fiberglass dune buggy was the 1964 "Meyers Manx" built by Bruce Meyers. Bruce Meyers designed his fiberglass bodies as a "kit car", using the Volkswagen Beetle chassis. Many other companies worldwide have been inspired by the Manx, making similar bodies and kits. These types of dune buggies are known as "clones". Sandrail A sandrail is a lightweight vehicle similar to a dune buggy, but designed specifically for operation on open sand. Sandrails are usually built as a spaceframe by welding steel tubes together. The name sandrail is due to the frame "rails" present. The advantage of this method is that the fabricator can change fundamental parts of the vehicle (usually the suspension and addition of a built-in roll cage). Sandrails, as per dune buggies, often have the engine located behind the driver. Sizes can vary from a small-engine one-seat size to four-seat vehicles with eight or more cylinders. A similar, more recent generation of off-road vehicle, often similar in appearance to a sandrail, but designed for a different use, is the "off road go-kart". The difference may be little more than fitting all-terrain tires instead of sand tires and the much smaller size of the engine. Military use Because of the advantages a buggy can afford on some terrain, they are also used by the military. The buggies built for the United States military used to be called Desert Patrol Vehicles (DPV) or Fast Attack Vehicles (FAV), and with the latest improvements are known as Light Strike Vehicles (LSV). They are used by United States Navy SEALs, the SAS, and other forces. Among the dune buggies used by the United States military is the Chenowth Advanced Light Strike Vehicle. The US Border Patrol also uses this (although it is not a military organization). In the United Kingdom, the SAS have used cut-down, light-weight all terrain vehicles for secret special operations "behind the lines" since early in World War II. A buggy was used by the British Special Air Service (SAS) forces during the Gulf War. A long-range special desert operations vehicle was developed in 1992 and nicknamed "pink panthers" because of their color, but these were only modified Land Rovers.	Technology	Motorized road transport	null
68513	https://en.wikipedia.org/wiki/Surface%20science	Surface science	Surface science is the study of physical and chemical phenomena that occur at the interface of two phases, including solid–liquid interfaces, solid–gas interfaces, solid–vacuum interfaces, and liquid–gas interfaces. It includes the fields of surface chemistry and surface physics. Some related practical applications are classed as surface engineering. The science encompasses concepts such as heterogeneous catalysis, semiconductor device fabrication, fuel cells, self-assembled monolayers, and adhesives. Surface science is closely related to interface and colloid science. Interfacial chemistry and physics are common subjects for both. The methods are different. In addition, interface and colloid science studies macroscopic phenomena that occur in heterogeneous systems due to peculiarities of interfaces. History The field of surface chemistry started with heterogeneous catalysis pioneered by Paul Sabatier on hydrogenation and Fritz Haber on the Haber process. Irving Langmuir was also one of the founders of this field, and the scientific journal on surface science, Langmuir, bears his name. The Langmuir adsorption equation is used to model monolayer adsorption where all surface adsorption sites have the same affinity for the adsorbing species and do not interact with each other. Gerhard Ertl in 1974 described for the first time the adsorption of hydrogen on a palladium surface using a novel technique called LEED. Similar studies with platinum, nickel, and iron followed. Most recent developments in surface sciences include the 2007 Nobel prize of Chemistry winner Gerhard Ertl's advancements in surface chemistry, specifically his investigation of the interaction between carbon monoxide molecules and platinum surfaces. Chemistry Surface chemistry can be roughly defined as the study of chemical reactions at interfaces. It is closely related to surface engineering, which aims at modifying the chemical composition of a surface by incorporation of selected elements or functional groups that produce various desired effects or improvements in the properties of the surface or interface. Surface science is of particular importance to the fields of heterogeneous catalysis, electrochemistry, and geochemistry. Catalysis The adhesion of gas or liquid molecules to the surface is known as adsorption. This can be due to either chemisorption or physisorption, and the strength of molecular adsorption to a catalyst surface is critically important to the catalyst's performance (see Sabatier principle). However, it is difficult to study these phenomena in real catalyst particles, which have complex structures. Instead, well-defined single crystal surfaces of catalytically active materials such as platinum are often used as model catalysts. Multi-component materials systems are used to study interactions between catalytically active metal particles and supporting oxides; these are produced by growing ultra-thin films or particles on a single crystal surface. Relationships between the composition, structure, and chemical behavior of these surfaces are studied using ultra-high vacuum techniques, including adsorption and temperature-programmed desorption of molecules, scanning tunneling microscopy, low energy electron diffraction, and Auger electron spectroscopy. Results can be fed into chemical models or used toward the rational design of new catalysts. Reaction mechanisms can also be clarified due to the atomic-scale precision of surface science measurements. Electrochemistry Electrochemistry is the study of processes driven through an applied potential at a solid–liquid or liquid–liquid interface. The behavior of an electrode–electrolyte interface is affected by the distribution of ions in the liquid phase next to the interface forming the electrical double layer. Adsorption and desorption events can be studied at atomically flat single-crystal surfaces as a function of applied potential, time and solution conditions using spectroscopy, scanning probe microscopy and surface X-ray scattering. These studies link traditional electrochemical techniques such as cyclic voltammetry to direct observations of interfacial processes. Geochemistry Geological phenomena such as iron cycling and soil contamination are controlled by the interfaces between minerals and their environment. The atomic-scale structure and chemical properties of mineral–solution interfaces are studied using in situ synchrotron X-ray techniques such as X-ray reflectivity, X-ray standing waves, and X-ray absorption spectroscopy as well as scanning probe microscopy. For example, studies of heavy metal or actinide adsorption onto mineral surfaces reveal molecular-scale details of adsorption, enabling more accurate predictions of how these contaminants travel through soils or disrupt natural dissolution–precipitation cycles. Physics Surface physics can be roughly defined as the study of physical interactions that occur at interfaces. It overlaps with surface chemistry. Some of the topics investigated in surface physics include friction, surface states, surface diffusion, surface reconstruction, surface phonons and plasmons, epitaxy, the emission and tunneling of electrons, spintronics, and the self-assembly of nanostructures on surfaces. Techniques to investigate processes at surfaces include surface X-ray scattering, scanning probe microscopy, surface-enhanced Raman spectroscopy and X-ray photoelectron spectroscopy. Analysis techniques The study and analysis of surfaces involves both physical and chemical analysis techniques. Several modern methods probe the topmost 1–10 nm of surfaces exposed to vacuum. These include angle-resolved photoemission spectroscopy (ARPES), X-ray photoelectron spectroscopy (XPS), Auger electron spectroscopy (AES), low-energy electron diffraction (LEED), electron energy loss spectroscopy (EELS), thermal desorption spectroscopy (TPD), ion scattering spectroscopy (ISS), secondary ion mass spectrometry, dual-polarization interferometry, and other surface analysis methods included in the list of materials analysis methods. Many of these techniques require vacuum as they rely on the detection of electrons or ions emitted from the surface under study. Moreover, in general ultra-high vacuum, in the range of 10−7 pascal pressure or better, it is necessary to reduce surface contamination by residual gas, by reducing the number of molecules reaching the sample over a given time period. At 0.1 mPa (10−6 torr) partial pressure of a contaminant and standard temperature, it only takes on the order of 1 second to cover a surface with a one-to-one monolayer of contaminant to surface atoms, so much lower pressures are needed for measurements. This is found by an order of magnitude estimate for the (number) specific surface area of materials and the impingement rate formula from the kinetic theory of gases. Purely optical techniques can be used to study interfaces under a wide variety of conditions. Reflection-absorption infrared, dual polarisation interferometry, surface-enhanced Raman spectroscopy and sum frequency generation spectroscopy can be used to probe solid–vacuum as well as solid–gas, solid–liquid, and liquid–gas surfaces. Multi-parametric surface plasmon resonance works in solid–gas, solid–liquid, liquid–gas surfaces and can detect even sub-nanometer layers. It probes the interaction kinetics as well as dynamic structural changes such as liposome collapse or swelling of layers in different pH. Dual-polarization interferometry is used to quantify the order and disruption in birefringent thin films. This has been used, for example, to study the formation of lipid bilayers and their interaction with membrane proteins. Acoustic techniques, such as quartz crystal microbalance with dissipation monitoring, is used for time-resolved measurements of solid–vacuum, solid–gas and solid–liquid interfaces. The method allows for analysis of molecule–surface interactions as well as structural changes and viscoelastic properties of the adlayer. X-ray scattering and spectroscopy techniques are also used to characterize surfaces and interfaces. While some of these measurements can be performed using laboratory X-ray sources, many require the high intensity and energy tunability of synchrotron radiation. X-ray crystal truncation rods (CTR) and X-ray standing wave (XSW) measurements probe changes in surface and adsorbate structures with sub-Ångström resolution. Surface-extended X-ray absorption fine structure (SEXAFS) measurements reveal the coordination structure and chemical state of adsorbates. Grazing-incidence small angle X-ray scattering (GISAXS) yields the size, shape, and orientation of nanoparticles on surfaces. The crystal structure and texture of thin films can be investigated using grazing-incidence X-ray diffraction (GIXD, GIXRD). X-ray photoelectron spectroscopy (XPS) is a standard tool for measuring the chemical states of surface species and for detecting the presence of surface contamination. Surface sensitivity is achieved by detecting photoelectrons with kinetic energies of about 10–1000 eV, which have corresponding inelastic mean free paths of only a few nanometers. This technique has been extended to operate at near-ambient pressures (ambient pressure XPS, AP-XPS) to probe more realistic gas–solid and liquid–solid interfaces. Performing XPS with hard X-rays at synchrotron light sources yields photoelectrons with kinetic energies of several keV (hard X-ray photoelectron spectroscopy, HAXPES), enabling access to chemical information from buried interfaces. Modern physical analysis methods include scanning-tunneling microscopy (STM) and a family of methods descended from it, including atomic force microscopy (AFM). These microscopies have considerably increased the ability of surface scientists to measure the physical structure of many surfaces. For example, they make it possible to follow reactions at the solid–gas interface in real space, if those proceed on a time scale accessible by the instrument.	Physical sciences	Subdisciplines	Chemistry
68518	https://en.wikipedia.org/wiki/Chemisorption	Chemisorption	Chemisorption is a kind of adsorption which involves a chemical reaction between the surface and the adsorbate. New chemical bonds are generated at the adsorbent surface. Examples include macroscopic phenomena that can be very obvious, like corrosion, and subtler effects associated with heterogeneous catalysis, where the catalyst and reactants are in different phases. The strong interaction between the adsorbate and the substrate surface creates new types of electronic bonds. In contrast with chemisorption is physisorption, which leaves the chemical species of the adsorbate and surface intact. It is conventionally accepted that the energetic threshold separating the binding energy of "physisorption" from that of "chemisorption" is about 0.5 eV per adsorbed species. Due to specificity, the nature of chemisorption can greatly differ, depending on the chemical identity and the surface structural properties. The bond between the adsorbate and adsorbent in chemisorption is either ionic or covalent. Uses An important example of chemisorption is in heterogeneous catalysis which involves molecules reacting with each other via the formation of chemisorbed intermediates. After the chemisorbed species combine (by forming bonds with each other) the product desorbs from the surface. Self-assembled monolayers Self-assembled monolayers (SAMs) are formed by chemisorbing reactive reagents with metal surfaces. A famous example involves thiols (RS-H) adsorbing onto the surface of gold. This process forms strong Au-SR bonds and releases H2. The densely packed SR groups protect the surface. Gas-surface chemisorption Adsorption kinetics As an instance of adsorption, chemisorption follows the adsorption process. The first stage is for the adsorbate particle to come into contact with the surface. The particle needs to be trapped onto the surface by not possessing enough energy to leave the gas-surface potential well. If it elastically collides with the surface, then it would return to the bulk gas. If it loses enough momentum through an inelastic collision, then it "sticks" onto the surface, forming a precursor state bonded to the surface by weak forces, similar to physisorption. The particle diffuses on the surface until it finds a deep chemisorption potential well. Then it reacts with the surface or simply desorbs after enough energy and time. The reaction with the surface is dependent on the chemical species involved. Applying the Gibbs energy equation for reactions: General thermodynamics states that for spontaneous reactions at constant temperature and pressure, the change in free energy should be negative. Since a free particle is restrained to a surface, and unless the surface atom is highly mobile, entropy is lowered. This means that the enthalpy term must be negative, implying an exothermic reaction. Physisorption is given as a Lennard-Jones potential and chemisorption is given as a Morse potential. There exists a point of crossover between the physisorption and chemisorption, meaning a point of transfer. It can occur above or below the zero-energy line (with a difference in the Morse potential, a), representing an activation energy requirement or lack of. Most simple gases on clean metal surfaces lack the activation energy requirement. Modeling For experimental setups of chemisorption, the amount of adsorption of a particular system is quantified by a sticking probability value. However, chemisorption is very difficult to theorize. A multidimensional potential energy surface (PES) derived from effective medium theory is used to describe the effect of the surface on absorption, but only certain parts of it are used depending on what is to be studied. A simple example of a PES, which takes the total of the energy as a function of location: where is the energy eigenvalue of the Schrödinger equation for the electronic degrees of freedom and is the ion interactions. This expression is without translational energy, rotational energy, vibrational excitations, and other such considerations. There exist several models to describe surface reactions: the Langmuir–Hinshelwood mechanism in which both reacting species are adsorbed, and the Eley–Rideal mechanism in which one is adsorbed and the other reacts with it. Real systems have many irregularities, making theoretical calculations more difficult: Solid surfaces are not necessarily at equilibrium. They may be perturbed and irregular, defects and such. Distribution of adsorption energies and odd adsorption sites. Bonds formed between the adsorbates. Compared to physisorption where adsorbates are simply sitting on the surface, the adsorbates can change the surface, along with its structure. The structure can go through relaxation, where the first few layers change interplanar distances without changing the surface structure, or reconstruction where the surface structure is changed. A direct transition from physisorption to chemisorption has been observed by attaching a CO molecule to the tip of an atomic force microscope and measuring its interaction with a single iron atom. For example, oxygen can form very strong bonds (~4 eV) with metals, such as Cu(110). This comes with the breaking apart of surface bonds in forming surface-adsorbate bonds. A large restructuring occurs by missing row. Dissociative chemisorption A particular brand of gas-surface chemisorption is the dissociation of diatomic gas molecules, such as hydrogen, oxygen, and nitrogen. One model used to describe the process is precursor-mediation. The absorbed molecule is adsorbed onto a surface into a precursor state. The molecule then diffuses across the surface to the chemisorption sites. They break the molecular bond in favor of new bonds to the surface. The energy to overcome the activation potential of dissociation usually comes from translational energy and vibrational energy. An example is the hydrogen and copper system, one that has been studied many times over. It has a large activation energy of 0.35 – 0.85 eV. The vibrational excitation of the hydrogen molecule promotes dissociation on low index surfaces of copper.	Physical sciences	Other separations	Chemistry
68520	https://en.wikipedia.org/wiki/Physisorption	Physisorption	Physisorption, also called physical adsorption, is a process in which the electronic structure of the atom or molecule is barely perturbed upon adsorption. Overview The fundamental interacting force of physisorption is Van der Waals force. Even though the interaction energy is very weak (~10–100 meV), physisorption plays an important role in nature. For instance, the van der Waals attraction between surfaces and foot-hairs of geckos (see Synthetic setae) provides the remarkable ability to climb up vertical walls. Van der Waals forces originate from the interactions between induced, permanent or transient electric dipoles. In comparison with chemisorption, in which the electronic structure of bonding atoms or molecules is changed and covalent or ionic bonds form, physisorption does not result in changes to the chemical bonding structure. In practice, the categorisation of a particular adsorption as physisorption or chemisorption depends principally on the binding energy of the adsorbate to the substrate, with physisorption being far weaker on a per-atom basis than any type of connection involving a chemical bond. Modeling by image charge To give a simple illustration of physisorption, we can first consider an adsorbed hydrogen atom in front of a perfect conductor, as shown in Fig. 1. A nucleus with positive charge is located at R = (0, 0, Z), and the position coordinate of its electron, r = (x, y, z) is given with respect to the nucleus. The adsorption process can be viewed as the interaction between this hydrogen atom and its image charges of both the nucleus and electron in the conductor. As a result, the total electrostatic energy is the sum of attraction and repulsion terms: The first term is the attractive interaction of the nucleus and its image charge, and the second term is due to the interaction of the electron and its image charge. The repulsive interaction is shown in the third and fourth terms arising from the interaction between the nucleus and the image electron, and, the interaction between the electron and the image nucleus, respectively. By Taylor expansion in powers of \|r\| / \|R\|, this interaction energy can be further expressed as: One can find from the first non-vanishing term that the physisorption potential depends on the distance Z between adsorbed atom and surface as Z−3, in contrast with the r−6 dependence of the molecular van der Waals potential, where r is the distance between two dipoles. Modeling by quantum-mechanical oscillator The van der Waals binding energy can be analyzed by another simple physical picture: modeling the motion of an electron around its nucleus by a three-dimensional simple harmonic oscillator with a potential energy Va: where me and ω are the mass and vibrational frequency of the electron, respectively. As this atom approaches the surface of a metal and forms adsorption, this potential energy Va will be modified due to the image charges by additional potential terms which are quadratic in the displacements: (from the Taylor expansion above.) Assuming the potential is well approximated as , where If one assumes that the electron is in the ground state, then the van der Waals binding energy is essentially the change of the zero-point energy: This expression also shows the nature of the Z−3 dependence of the van der Waals interaction. Furthermore, by introducing the atomic polarizability, the van der Waals potential can be further simplified: where is the van der Waals constant which is related to the atomic polarizability. Also, by expressing the fourth-order correction in the Taylor expansion above as (aCvZ0) / (Z4), where a is some constant, we can define Z0 as the position of the dynamical image plane and obtain The origin of Z0 comes from the spilling of the electron wavefunction out of the surface. As a result, the position of the image plane representing the reference for the space coordinate is different from the substrate surface itself and modified by Z0. Table 1 shows the jellium model calculation for van der Waals constant Cv and dynamical image plane Z0 of rare gas atoms on various metal surfaces. The increasing of Cv from He to Xe for all metal substrates is caused by the larger atomic polarizability of the heavier rare gas atoms. For the position of the dynamical image plane, it decreases with increasing dielectric function and is typically on the order of 0.2 Å. Physisorption potential Even though the van der Waals interaction is attractive, as the adsorbed atom moves closer to the surface the wavefunction of electron starts to overlap with that of the surface atoms. Further the energy of the system will increase due to the orthogonality of wavefunctions of the approaching atom and surface atoms. This Pauli exclusion and repulsion are particularly strong for atoms with closed valence shells that dominate the surface interaction. As a result, the minimum energy of physisorption must be found by the balance between the long-range van der Waals attraction and short-range Pauli repulsion. For instance, by separating the total interaction of physisorption into two contributions—a short-range term depicted by Hartree–Fock theory and a long-range van der Waals attraction—the equilibrium position of physisorption for rare gases adsorbed on jellium substrate can be determined. Fig. 2 shows the physisorption potential energy of He adsorbed on Ag, Cu, and Au substrates which are described by the jellium model with different densities of smear-out background positive charges. It can be found that the weak van der Waals interaction leads to shallow attractive energy wells (<10 meV). One of the experimental methods for exploring physisorption potential energy is the scattering process, for instance, inert gas atoms scattered from metal surfaces. Certain specific features of the interaction potential between scattered atoms and surface can be extracted by analyzing the experimentally determined angular distribution and cross sections of the scattered particles. Quantum mechanical – thermodynamic modelling for surface area and porosity Since 1980 two theories were worked on to explain adsorption and obtain equations that work. These two are referred to as the chi hypothesis, the quantum mechanical derivation, and excess surface work, ESW. Both these theories yield the same equation for flat surfaces: Where U is the unit step function. The definitions of the other symbols is as follows: where "ads" stands for "adsorbed", "m" stands for "monolayer equivalence" and "vap" is reference to the vapor pressure ("ads" and "vap" are the latest IUPAC convention but "m" has no IUAPC equivalent notation) of the liquid adsorptive at the same temperature as the solid sample. The unit function creates the definition of the molar energy of adsorption for the first adsorbed molecule by: The plot of adsorbed versus is referred to as the chi plot. For flat surfaces, the slope of the chi plot yields the surface area. Empirically, this plot was notice as being a very good fit to the isotherm by Polanyi and also by deBoer and Zwikker but not pursued. This was due to criticism in the former case by Einstein and in the latter case by Brunauer. This flat surface equation may be used as a "standard curve" in the normal tradition of comparison curves, with the exception that the porous sample's early portion of the plot of versus acts as a self-standard. Ultramicroporous, microporous and mesoporous conditions may be analyzed using this technique. Typical standard deviations for full isotherm fits including porous samples are typically less than 2%. A typical fit to good data on a homogeneous non-porous surface is shown in figure 3. The data is by Payne, Sing and Turk and was used to create the -s standard curve. Unlike the BET, which can only be at best fit over the range of 0.05 to 0.35 of P/Pvap, the range of the fit is the full isotherm. Comparison with chemisorption Physisorption is a general phenomenon and occurs in any solid/fluid or solid/gas system. Chemisorption is characterized by chemical specificity. In physisorption, perturbation of the electronic states of adsorbent and adsorbate is minimal. The adsorption forces include London Forces, dipole-dipole attractions, dipole-induced attraction and "hydrogen bonding." For chemisorption, changes in the electronic states may be detectable by suitable physical means, in other words, chemical bonding. Typical binding energy of physisorption is about 10–300 meV and non-localized. Chemisorption usually forms bonding with energy of 1–10 eV and localized. The elementary step in physisorption from a gas phase does not involve activation energy. Chemisorption often involves an activation energy. For physisorption gas phase molecules, adsorbates, form multilayer adsorption unless physical barriers, such as porosity, interfere. In chemisorption, molecules are adsorbed on the surface by valence bonds and only form monolayer adsorption. A direct transition from physisorption to chemisorption has been observed by attaching a CO molecule to the tip of an atomic force microscope and measuring its interaction with a single iron atom. This effect was observed in the late 1960s for benzene from field emission as reported by Condon and ESR measurements as reported by Moyes and Wells. Another way of looking at this is that chemisorption alters the topology of the electrons in the adsorbate molecule (by the process of chemical reaction) but physisorption does not.	Physical sciences	Other separations	Chemistry
68630	https://en.wikipedia.org/wiki/Alarm%20device	Alarm device	An alarm device is a mechanism that gives an audible, visual, combination, or other kind of alarm signal to alert someone to a problem or condition that requires urgent attention. Etymology The word alarm comes from the Old French a l'arme meaning "to the arms", or "to the weapons", telling armed men to pick up their weapons and get ready for action because an enemy may have suddenly appeared. The word alarum is an archaic form of alarm. It was sometimes used as a call to arms in the stage directions of Elizabethan dramas. The term comes from the Italian all'armi and appears 89 times in Shakespeare's First Folio. Often explained as the off-stage sounds of conflict or disturbance, recent research suggests a bell or drum may have been used to rouse soldiers from sleep. History and development Early alarm devices were often bells, drums, other musical instruments, or any items which made unusual loud noises that attracted the attention of the surrounding population. Whistles were used by police in the 19th century. Steam whistles have been used on locomotives, ships, and in factories as alarm devices. With the advent of electricity, a variety of other alerting devices have been invented, such as buzzers, klaxons, sirens, horns, flashing and coloured lights, and other all-purpose alarms. Alarm devices can be fitted to buildings as well as vehicles. Many buildings are fitted with fire alarms, ranging from a self-contained domestic smoke detector to a sophisticated alarm system that can operate building fire fighting systems automatically to extinguish fires with water or inert gases. Many industries have developed standards for alarm devices, and the colours red, blue and amber are generally recognized as alarm device-related colours, with flashing lights often indicating urgent conditions. Responses to an alarm Human reactions to an alarm will often depend on upbringing, psychological training, or the behavior of others in the environment. Consequently, the ability to test an alarm and hold regular drills to practice an appropriate response may be provided as part of an alarm system. Alarm devices that are intended to cause the evacuation of an occupied building, such as fire alarms, may be deliberately designed to make remaining in the space difficult or even painful in order to encourage occupants to leave. Some alarms may startle and cause a fight-or-flight response in humans; a person under this mindset will panic and either flee the perceived danger or attempt to eliminate it, often ignoring rational thought in either case. A person in such a state can be characterized as "alarmed". False alarms With any kind of alarm, you must balance between the danger of false alarms (called "false positives") — the signal going off in the absence of a problem — or an alarm failing to signal an actual problem (called a "false negative"). False alarms can waste resources expensively and can even be dangerous. For example, false alarms of a fire can waste firefighter manpower, making them unavailable for a real fire, and risk injury to firefighters and others as the fire engines race to the alleged fire's location. In addition, false alarms may acclimatise people to ignore alarm signals, and thus possibly to ignore an actual emergency: Aesop's fable of The Boy Who Cried Wolf exemplifies this problem. A false alarm is one of the most significant issues with conventional alarm systems. They can be triggered for several reasons, such as the movement of pets, typing in the wrong security codes, or loud sounds from windows or doors. In the case of fire alarms, aerosol sprays, smoking, or burning food can all lead to a false alarm. Many avoid the risk of false alarms by ensuring their alarms are secured in an appropriate location, such as placing a smoke detector or fire alarm away from the kitchen where smoke from burned food or large quantities of steam which may trigger a false alarm are common occurrences. In the case of a security alarm, an additional monitoring station which assesses whether there is a legitimate need for help can reduce false alarms. Devices There are many kinds of alarm devices. The most common types include: an alarm clock that sounds an alarm at a pre-set time, often used to wake a person up or remind them of an event. a fire alarm which is used to give occupants of a building early warning of a potential fire and give them time to evacuate. warning devices on a vehicle that sound when it is moving in an unexpected direction, such as reversing, a siren, often accompanied by flashing coloured lights, on emergency vehicles responding to an emergency Alarm devices, by category, include: burglar alarms, designed to warn of burglaries. This is often a silent alarm; law enforcement or guards are warned without alerting the burglar, which increases the chances of stopping the theft while in progress. alarm clocks can beep, buzz or ring at a set time to wake a person up or for other reminders distributed control systems (DCS), found in nuclear power plants, refineries and chemical facilities, also generate alarms to direct the operator's attention to an important event that they need to address. alarms in an operation and maintenance (O&M) monitoring system, which alerts an operator to a malfunction of a particular part of the system under monitoring. first-out alarm safety alarms, which go off if a dangerous condition occurs. Common public safety alarms include: civil defense siren, also known as tornado sirens or air raid sirens fire alarm systems fire alarm notification appliance "Multiple-alarm fire", a locally specific measure of the severity of a fire and the fire-department reaction required. smoke detector car alarms autodialer alarm, also known as community alarm personal alarm Video alarm verification systems provides instant notifications upon the detection of a possible threat verified through a video feed. tocsin – a historical alarm mechanism	Technology	Mechanisms	null
68659	https://en.wikipedia.org/wiki/Ammonoidea	Ammonoidea	Ammonoids are extinct, (typically) coiled-shelled cephalopods comprising the subclass Ammonoidea. They are more closely related to living octopuses, squid, and cuttlefish (which comprise the clade Coleoidea) than they are to nautiluses (family Nautilidae). The earliest ammonoids appeared during the Emsian stage of the Early Devonian, with the last species vanishing during or soon after the Cretaceous–Paleogene extinction event. They are often called ammonites, which is most frequently used for members of the order Ammonitida, the only remaining group of ammonoids from the Jurassic up until their extinction. Ammonoids exhibited considerable diversity over their evolutionary history, with over 10,000 species having been described. Ammonoids are excellent index fossils, and they have been frequently used to link rock layers in which a particular species or genus is found to specific geologic time periods. Their fossil shells usually take the form of planispirals, although some helically spiraled and nonspiraled forms (known as heteromorphs) have been found, primarily during the Cretaceous period. The name "ammonite", from which the scientific term is derived, was inspired by the spiral shape of their fossilized shells, which somewhat resemble tightly coiled rams' horns. Pliny the Elder ( 79 AD near Pompeii) called fossils of these animals ("horns of Ammon") because the Egyptian god Ammon (Amun) was typically depicted wearing rams' horns. Often, the name of an ammonite genus ends in -ceras, which is from () meaning "horn". Classification Orders and suborders The Ammonoidea can be divided into six orders, listed here starting with the most primitive and going to the more derived: Agoniatitida, Lower Devonian – Middle Devonian Clymeniida, Upper Devonian Goniatitida, Middle Devonian – Upper Permian Prolecanitida, Upper Devonian – Upper Triassic Ceratitida, Upper Permian – Upper Triassic Ammonitida, Lower Jurassic – Lower Paleocene In some classifications, these are left as suborders, included in only three orders: Goniatitida, Ceratitida and Ammonitida. The classification of ammonoids is based in part on the ornamentation and structure of the septa comprising their shells' gas chambers. Taxonomy of the Treatise on Invertebrate Paleontology The Treatise on Invertebrate Paleontology (Part L, 1957) divides the Ammonoidea, regarded simply as an order, into eight suborders, the Anarcestina, Clymeniina, Goniatitina and Prolecanitina from the Paleozoic; the Ceratitina from the Triassic; and the Ammonitina, Lytoceratina and Phylloceratina from the Jurassic and Cretaceous. In subsequent taxonomies, these are sometimes regarded as orders within the subclass Ammonoidea. Shell anatomy and diversity Basic shell anatomy The chambered part of the ammonite shell is called a phragmocone. It contains a series of progressively larger chambers, called camerae (sing. camera) that are divided by thin walls called septa (sing. septum). Only the last and largest chamber, the body chamber, was occupied by the living animal at any given moment. As it grew, it added newer and larger chambers to the open end of the coil. Where the outer whorl of an ammonite shell largely covers the preceding whorls, the specimen is said to be involute (e.g., Anahoplites). Where it does not cover those preceding, the specimen is said to be evolute (e.g., Dactylioceras). A thin living tube called a siphuncle passed through the septa, extending from the ammonite's body into the empty shell chambers. Through a hyperosmotic active transport process, the ammonite emptied water out of these shell chambers. This enabled it to control the buoyancy of the shell and thereby rise or descend in the water column. A primary difference between ammonites and nautiloids is the siphuncle of ammonites (excepting Clymeniina) runs along the ventral periphery of the septa and camerae (i.e., the inner surface of the outer axis of the shell), while the siphuncle of nautiloids runs more or less through the center of the septa and camerae. Septa and suture patterns Ammonites (subclass Ammonoidea) can be distinguished by their septa, the dividing walls that separate the chambers in the phragmocone, by the nature of their sutures where the septa join the outer shell wall, and in general by their siphuncles. Ammonoid septa characteristically have bulges and indentations and are to varying degrees convex when seen from the front, distinguishing them from nautiloid septa, which are typically simple concave, dish-shaped structures. The topology of the septa, especially around the rim, results in the various suture patterns found. The septal curvature in nautiloids and ammonoids also differ in that the septa curves towards the opening in nautiloids, and away from the opening in ammоnoids. While nearly all nautiloids show gently curving sutures, the ammonoid suture line (the intersection of the septum with the outer shell) is variably folded, forming saddles ("peaks" that point towards the aperture) and lobes ("valleys" which point away from the aperture). The suture line has four main regions. The external or ventral region refers to sutures along the lower (outer) edge of the shell, where the left and right suture lines meet. The external (or ventral) saddle, when present, lies directly on the lower midline of the shell. As a result, it is often called the median saddle. On suture diagrams the median saddle is supplied with an arrow which points towards the aperture. The median saddle is edged by fairly small external (or ventral) lobes. The earliest ammonoids lacked a median saddle and instead had a single midline ventral lobe, which in later forms is split into two or more components. The lateral region involves the first saddle and lobe pair past the external region as the suture line extends up the side of the shell. The lateral saddle and lobe are usually larger than the ventral saddle and lobe. Additional lobes developing towards the inner edge of a whorl are labelled umbilical lobes, which increase in number through ammonoid evolution as well as an individual ammonoid's development. In many cases the distinction between the lateral and umbilical regions are unclear; new umbilical features can develop from subdivisions of other umbilical features, or from subdivisions of lateral features. Lobes and saddles which are so far towards the center of the whorl that they are covered up by succeeding whorls are labelled internal (or dorsal) lobes and saddles. Three major types of suture patterns are found in the Ammonoidea: Goniatitic – numerous undivided lobes and saddles. This pattern is characteristic of the Paleozoic ammonoids (orders Agoniatitida, Clymeniida, Goniatitida, and Prolecanitida). Ceratitic – lobes have subdivided tips, giving them a saw-toothed appearance. The saddles are rounded and undivided. This suture pattern is characteristic of Triassic ammonoids in the order Ceratitida. It appears again in the Cretaceous "pseudoceratites". Ammonitic – lobes and saddles are much subdivided (fluted); subdivisions are usually rounded instead of saw-toothed. Ammonoids of this type are the most important species from a biostratigraphical point of view. This suture type is characteristic of Jurassic and Cretaceous ammonoids, but extends back all the way to the Permian. Siphuncle The siphuncle in most ammonoids is a narrow tubular structure that runs along the shell's outer rim, known as the venter, connecting the chambers of the phragmocone to the body or living chamber. This distinguishes them from living nautiloides (Nautilus and Allonautilus) and typical Nautilida, in which the siphuncle runs through the center of each chamber. However the very earliest nautiloids from the Late Cambrian and Ordovician typically had ventral siphuncles like ammonites, although often proportionally larger and more internally structured. The word "siphuncle" comes from the Neo-Latin siphunculus, meaning "little siphon". Sexual dimorphism One feature found in shells of the modern Nautilus is the variation in the shape and size of the shell according to the sex of the animal, the shell of the male being slightly smaller and wider than that of the female. This sexual dimorphism is thought to be an explanation for the variation in size of certain ammonite shells of the same species, the larger shell (the macroconch) being female, and the smaller shell (the microconch) being male. This is thought to be because the female required a larger body size for egg production. A good example of this sexual variation is found in Bifericeras from the early part of the Jurassic period of Europe. Only recently has sexual variation in the shells of ammonites been recognized. The macroconch and microconch of one species were often previously mistaken for two closely related but different species occurring in the same rocks. However, because the dimorphic sizes are so consistently found together, they are more likely an example of sexual dimorphism within the same species. Whorl width in the body chamber of many groups of ammonites, as expressed by the width:diameter ratio, is another sign of dimorphism. This character has been used to separate "male" (Largiventer conch "L") from "female" (Leviventer conch "l"). Variations in shape The majority of ammonite species feature planispiral shells, tightly coiled in a flat plane. The most fundamental difference in spiral form is how strongly successive whorls expand and overlap their predecessors. This can be inferred by the size of the umbilicus, the sunken-in inner part of the coil, exposing older and smaller whorls. Evolute shells have very little overlap, a large umbilicus, and many exposed whorls. Involute shells have strong overlap, a small umbilicus, and only the largest and most recent whorls are exposed. Shell structure can be broken down further by the width of the shell, with implications for hydrodynamic efficiency. Major shell forms include: Oxycone – Strongly involute and very narrow, with sharp ventral keels and a streamlined, lenticular (lens-shaped) cross-section. These ammonoids are estimated to be nektonic (well-adapted to rapid active swimming), as their shell form incurs very little drag and allows for efficient, stable coasting even in turbulent flow regimes. Serpenticone – Strongly evolute and fairly narrow (discoidal) in width. Historically assumed to be primarily planktonic (free-floating drifters), a nektonic lifestyle is also plausible for many species. Thanks to their flattened shape, these ammonoids accelerate effectively, though their large umbilicus introduces more drag in successive thrusts. Relative to oxycones, serpenticones take less effort to rotate around the transverse axis (pitch). Serpenticone ammonites resemble coiled snakes and are abundant in the Jurassic rocks of Europe. Carved serpenticones fulfill the role of the "snakestones" in medieval folklore. Spherocone – Moderately involute and quite broad, globular (nearly spherical) in overall shape. Their semi-spherical shape is the most efficient for moving in laminar water (with a low Reynolds number) or migrating vertically through the water column. Though less hydrodynamically stable than other forms, this may be advantageous in certain situations, as spherocones can easily rotate around both the transverse axis and the vertical axis (yaw). Platycone – Intermediate between serpenticones and oxycones: narrow and moderately involute. Discocone – Intermediate between oxycones and spherocones: involute and moderately broad. The modern Nautilus is an example of a discocone cephalopod. Planorbicone – Intermediate between serpenticones and spherocones: Moderately broad, evolute to involute. Wider and more involute ammonoids on the serpenticone-spherocone spectrum are termed Cadicones. Ammonites vary greatly in the ornamentation (surface relief) of their shells. Some may be smooth and relatively featureless, except for growth lines, resembling that of the modern Nautilus. In others, various patterns of spiral ridges, ribs, nodes, or spines are presented. This type of complex ornamentation of the shell is especially evident in the later ammonites of the Cretaceous. Heteromorphs Ammonoids with a shell shape diverging from the typical planispiral form are known as heteromorphs, instead forming a conch with detached whorls (open coiling) or non-planispiral coiling. These types of shells evolved four times in ammonoids, with the first forms appearing already in the Devonian period. In late Norian age in Triassic the first heteromorph ammonoid fossils belongs to the genus Rhabdoceras. The three other heteromorphic genera were Hannaoceras, Cochloceras and Choristoceras. All of them went extinct at the end of Triassic. In the Jurassic an uncoiled shell was found in the Spiroceratoidea, but by the end of Cretaceous the only heteromorph ammonites remaining belonged to the suborder Ancyloceratina. One example is Baculites, which has a nearly straight shell convergent with the older orthocone nautiloids. Still other species' shells are coiled helically (in two dimensions), similar in appearance to some gastropods (e.g., Turrilites and Bostrychoceras). Some species' shells are even initially uncoiled, then partially coiled, and finally straight at maturity (as in Australiceras). Perhaps the most extreme and bizarre-looking example of a heteromorph is Nipponites, which appears to be a tangle of irregular whorls lacking any obvious symmetric coiling. Upon closer inspection, though, the shell proves to be a three-dimensional network of connected "U" shapes. occurs in rocks of the upper part of the Cretaceous in Japan and the United States. Aptychus Some ammonites have been found in association with a single horny plate or a pair of calcitic plates. In the past, these plates were assumed to serve in closing the opening of the shell in much the same way as an operculum, but more recently they are postulated to have been a jaw apparatus. The plates are collectively termed the aptychus or aptychi in the case of a pair of plates, and anaptychus in the case of a single plate. The paired aptychi were symmetric to one another and equal in size and appearance. Anaptychi are relatively rare as fossils. They are found representing ammonites from the Devonian period through those of the Cretaceous period. Calcified aptychi only occur in ammonites from the Mesozoic era. They are almost always found detached from the shell, and are only very rarely preserved in place. Still, sufficient numbers have been found closing the apertures of fossil ammonite shells as to leave no doubt as to their identity as part of the anatomy of an ammonite. Large numbers of detached aptychi occur in certain beds of rock (such as those from the Mesozoic in the Alps). These rocks are usually accumulated at great depths. The modern Nautilus lacks any calcitic plate for closing its shell, and only one extinct nautiloid genus is known to have borne anything similar. Nautilus does, however, have a leathery head shield (the hood) which it uses to cover the opening when it retreats inside. There are many forms of aptychus, varying in shape and the sculpture of the inner and outer surfaces, but because they are so rarely found in position within the shell of the ammonite it is often unclear to which species of ammonite one kind of aptychus belongs. A number of aptychi have been given their own genus and even species names independent of their unknown owners' genus and species, pending future discovery of verified occurrences within ammonite shells. Soft tissues, life appearance and ecology Because ammonites and their close relatives are extinct, little is known about their way of life. Their soft body parts are very rarely preserved in any detail. Nonetheless, much has been worked out by examining ammonoid shells and by using models of these shells in water tanks. Although ammonites do occur in exceptional lagerstatten such as the Solnhofen Limestone, their soft-part record is surprisingly sparse. Beyond a tentative ink sac and possible digestive organs, no soft parts were known until 2021. In this year an isolated specimen showing some of the internal soft anatomy including organs was described. When neutron imaging was used on a fossil found in 1998, part of the musculature became visible and showed they were able to retract themselves into the shell for protection, and that the retractor muscles and hyponome that work together to enable jet propulsion in nautilus worked independently in ammonites. The soft body of the creature occupied the largest segments of the shell at the end of the coil. The smaller earlier segments were walled off and the animal could maintain its buoyancy by filling them with gas. Thus, the smaller sections of the coil would have floated above the larger sections. The reproductive organs show possible traces of spermatophores, which would support the hypothesis that the microconchs were males. They likely bore a radula and beak, and marginal siphuncle. They operated by direct development with sexual reproduction, were carnivorous, and had a crop for food storage. They are unlikely to have dwelt in fresh or brackish water. Many ammonites were likely filter feeders, so adaptations associated with this lifestyle like sieves probably occurred.A 2021 study reported specimens of the scaphitid ammonite genera Rhaeboceras and Hoploscaphites with mineralised hooks, which were likely present on the ends of a pair of enlarged tentacles. However, these mineralised hooks appear to be present only in scaphitids and were not typical of ammonites as a whole. The number of arms has been subject considerable speculation, with different artists either opting for a nautilus-like restoration with many arms, or a more squid-like restoration with much fewer arms, with a 1996 study suggesting that they probably had 10 arms like modern squid, cuttlefish and octopuses, but that nothing could be said for certain. Paleontologist Mark Witton has stated that "The basic details of ammonite life appearance are far from clear . . . While we can be certain that a squid-like organism lived in the last chamber of their shells . . . little else can be said with certainty about their appearance. ... Despite being creatures which occur so commonly as fossils that it seems like we should know everything about them, ammonites are creatures fraught with uncertainty for artists and palaeontologists alike. Until new data comes to light, all life reconstructions of ammonites should be taken as extremely tentative, almost speculative renditions of their actual appearance." Many ammonoids probably lived in the open water of ancient seas, rather than at the sea bottom, because their fossils are often found in rocks laid down under conditions where no bottom-dwelling life is found. In general, they appear to have inhabited the upper of the water column. Many of them (such as Oxynoticeras) are thought to have been good swimmers, with flattened, discus-shaped, streamlined shells, although some ammonoids were less effective swimmers and were likely to have been slow-swimming bottom-dwellers. Synchrotron analysis of an aptychophoran ammonite revealed remains of isopod and mollusc larvae in its buccal cavity, indicating at least this kind of ammonite fed on plankton. They may have avoided predation by squirting ink, much like modern cephalopods; ink is occasionally preserved in fossil specimens. Many ammonite shells have been found with round holes once interpreted as a result of limpets attaching themselves to the shells. However, the triangular formation of the holes, their size and shape, and their presence on both sides of the shells, corresponding to the upper and lower jaws, is more likely evidence of the bite of a medium-sized mosasaur preying upon ammonites. Some ammonites appear to have lived in cold seeps and even reproduced there. Size The smallest ammonoid was Maximites from the Upper Carboniferous. Adult specimens reached only in shell diameter. Few of the ammonites occurring in the lower and middle part of the Jurassic period reached a size exceeding in diameter. Much larger forms are found in the later rocks of the upper part of the Jurassic and the lower part of the Cretaceous, such as Titanites from the Portland Stone of Jurassic of southern England, which is often in diameter, and Parapuzosia seppenradensis of the Cretaceous period of Germany, which is one of the largest-known ammonites, sometimes reaching in diameter. The largest-documented North American ammonite is Parapuzosia bradyi from the Cretaceous, with specimens measuring in diameter. Evolutionary history Ammonoids are widely thought to have originated from straight-shelled (orthocone) "nautiloids" belong to Bactritida during the early Devonian (Emsian), with transitional fossils showing the transition from a straight shell, to a curved (cyrtoconic) shell to a relaxed (gyroconic) spiral and finally to a tight spiral. The Kellwasser Event at the end of the Frasnian let to a dramatic decline in ammonoid diversity, with only a handful of lineages belong to Tornoceratina (a subgroup of Goniatites) surviving, becoming ancestral to all later ammonoids. Ammonoids rediversified during the following Famennian, which also saw the radical shift of the siphuncle from a lower (ventral) to upper (dorsal) position. Ammonites were nearly completely exterminated by the Hangenberg Event at the end of the Devonian, with only a handful of lineages surviving, with one of the surviving goniatite lineages becoming ancestral to all post-early Carboniferous and later ammonoids. Ammonoids again rediversified during the Early Carboniferous. During the Carboniferous ammonoids underwent alternating periods of diversification and decline, and during the late Carboniferous ammonioid diversity became concentrated in a few geographical regions. During the Permian, the Capitanian mass extinction event severely reduced the diversity of Goniatitida and Prolecanitida, while the Ceratitida, which originated during the Middle Permian, likely from the Daraelitidae, was largely unaffected and radiated in the Late Permian, becoming the dominant group of ammonoids in this period represented by two groups, the araxoceratids and xenodiscids. The end-Permian mass extinction again reduced ammonoids to the verge of extinction, though both main ceratitd lineages lineages survived, though the xenodiscids were more successful and ancestral to all later ammonoids. Ammonites were devastated by the end-Triassic extinction, with only a handful of genera belonging to the family Psiloceratidae of the suborder Phylloceratina surviving and becoming ancestral to all later Jurassic and Cretaceous ammonites. Ammonites explosively diversified during the Early Jurassic, with the orders Psiloceratina, Ammonitina, Lytoceratina, Haploceratina, Perisphinctina and Ancyloceratina all appearing during the Jurassic. Heteromorph ammonites (ammonites with open or non-spiral coiling) of the order Ancyloceratina became common during the Cretaceous period. At least 57 species of ammonites, which were widespread and belonged to six superfamilies, were extant during the last 500,000 years of the Cretaceous, indicating that ammonites remained highly diverse until the very end of their existence. All ammonites were wiped out during or shortly after the K-Pg extinction event, caused by the Chicxulub impact. It has been suggested that ocean acidification generated by the impact played a key role in their extinction, as the larvae of ammonites were likely small and planktonic, and would have been heavily affected. Nautiloids, exemplified by modern nautiluses, are conversely thought to have had a reproductive strategy in which eggs were laid in smaller batches many times during the lifespan, and on the sea floor well away from any direct effects of such a bolide strike, and thus survived. Many ammonite species were filter feeders, so they might have been particularly susceptible to marine faunal turnovers and climatic change. Some reports suggest that a few ammonite species may have persisted into the very early Danian stage of the Paleocene, before going extinct. Cultural significance In medieval Europe, fossilised ammonites were thought to be petrified coiled snakes, and were called "snakestones" or, more commonly in medieval England, "serpentstones". They were considered to be evidence for the actions of saints, such as Hilda of Whitby, a myth referenced in Sir Walter Scott's Marmion, and Saint Patrick, and were held to have healing or oracular powers. Traders would occasionally carve the head of a snake onto the empty, wide end of the ammonite fossil, and then sell them as petrified snakes. In other cases, the snake's head would be simply painted on. Others believed ammonites, which they referred to as "salagrana" were composed of fossilized worm dung, and could be used to ward off witches. Ammonites from the Gandaki River in Nepal are known as Shaligrams, and are believed by Hindus to be a concrete manifestation of Vishnu.	Biology and health sciences	Cephalopods	Animals
68674	https://en.wikipedia.org/wiki/Anise	Anise	Anise (; ), also called aniseed or rarely anix, is a flowering plant in the family Apiaceae native to the eastern Mediterranean region and Southwest Asia. The flavor and aroma of its seeds have similarities with some other spices and herbs, such as star anise, fennel, liquorice, and tarragon. It is widely cultivated and used to flavor food, candy, and alcoholic drinks, especially around the Mediterranean. Etymology The name "anise" is derived via Old French from the Latin words or from Greek ánēthon referring to dill. An obsolete English word for anise is anet, also coming from anīsum. Botany Anise is an herbaceous annual plant growing to or more. The leaves at the base of the plant are simple, long and shallowly lobed, while leaves higher on the stems are feathery or lacy, pinnate, divided into numerous small leaflets. Both leaves and flowers are produced in large, loose clusters. The flowers are either white or yellow, approximately in diameter, produced in dense umbels. The fruit is a dry oblong and curved schizocarp, long, usually called "aniseed". Ecology Anise is a food plant for the larvae of some Lepidoptera species (butterflies and moths), including the lime-speck pug and wormwood pug. Cultivation Anise was first cultivated in Egypt and the Middle East, and was brought to Europe for its medicinal value. It has been cultivated in Egypt for approximately 4,000 years. Anise plants grow best in light, fertile, well-drained soil. The seeds should be planted as soon as the ground warms up in spring. Because the plants have a taproot, they do not transplant well after being established so they should either be started in their final location or be transplanted while the seedlings are still small. Production Western cuisines have long used anise to flavor dishes, drinks, and candies. The word is used for both the species of herb and its licorice-like flavor. The most powerful flavor component of the essential oil of anise, anethole, is found in both anise and an unrelated spice indigenous to South China called star anise (Illicium verum) widely used in South Asian, Southeast Asian and East Asian dishes. Star anise is considerably less expensive to produce and has gradually displaced P. anisum in Western markets. While formerly produced in larger quantities, by 1999 world production of the essential oil of anise was only 8 tons, compared to 400 tons of star anise. Uses Composition As with all spices, the composition of anise varies considerably with origin and cultivation method. These are typical values for the main constituents. Moisture: 9–13% Protein: 18% Fatty oil: 8–23% Essential oil: 2–7% Starch: 5% N-free extract: 22–28% Crude fibre: 12–25% In particular, the anise seeds products should also contain more than 0.2 milliliter volatile oil per 100 grams of spice. Culinary Anise is sweet and aromatic, distinguished by its characteristic flavor. The seeds, whole or ground, are used for preparation of teas and tisanes (alone or in combination with other aromatic herbs), as well many regional and ethnic confectioneries, including black jelly beans (often marketed as licorice-flavored), British aniseed balls, aniseed twists and "troach" drops, Australian humbugs, New Zealand aniseed wheels, Italian pizzelle and biscotti, German Pfeffernüsse and Springerle, Austrian Anisbögen, Dutch muisjes, New Mexican bizcochitos and Peruvian picarones. The culinary uses of anise are not limited only to sweets and confections, as it is a key ingredient in Mexican atole de anís and champurrado, which is similar to hot chocolate. In India and Pakistan, it is taken as a digestive after meals, used in brines in the Italian region of Apulia and as a flavoring agent in Italian sausage, pepperoni and other Italian processed meat products. The freshly chopped leaves are added to cheese spreads, dips or salads, while roots and stems impart a mild licorice flavor to soups and stews. The ancient Romans often served spiced cakes with aniseed called at the end of feasts as a digestive. This tradition of serving cake at the end of festivities is the basis for the tradition of serving cake at weddings. Liquor Anise is used to flavor Greek and Bulgarian ; Italian ; French , , and ; Spanish , , , and Herbs de Majorca; Turkish and Armenian ; Lebanese, Egyptian, Syrian, Jordanian, Palestinian and Israeli ; and Algerian . Outside the Mediterranean region, it is found in Colombian and Mexican . These liquors are clear, but on addition of water become cloudy, a phenomenon known as the ouzo effect. Anise is used together with other herbs and spices in some root beers, such as Virgil's in the United States. Traditional medicine The main use of anise in traditional European herbal medicine was for its carminative effect (reducing flatulence), as noted by John Gerard in his Great Herball, an early encyclopedia of herbal medicine: The seed wasteth and consumeth winde, and is good against belchings and upbraidings of the stomach, alaieth gripings of the belly, provoketh urine gently, maketh abundance of milke, and stirreth up bodily lust: it staieth the laske (diarrhea), and also the white flux (leukorrhea) in women. According to Pliny the Elder, anise was used as a cure for sleeplessness, chewed with alexanders and a little honey in the morning to freshen the breath, and, when mixed with wine, as a remedy for asp bites (N.H. 20.72). In 19th-century medicine, anise was prepared as ("Water of Anise") in doses of an ounce or more and as ("Spirit of Anise") in doses of 5–20 minims. In Turkish folk medicine, its seeds have been used as an appetite stimulant, tranquilizer or diuretic. Essential oil Anise essential oil can be obtained from the fruits by either steam distillation or extraction using supercritical carbon dioxide. The yield of essential oil is influenced by the growing conditions and extraction process, with supercritical extraction being more efficient. Regardless of the method of isolation the main component of the oil is anethole (80–90%), with minor components including 4-anisaldehyde, estragole and pseudoisoeugenyl-2-methylbutyrates amongst others. (Alternately found by Orav et al. 2008 to be 2–6% extracted oil by weight of raw seed material, 74–94% being trans-anethole and the remaining fraction estragole (methylchavicol), anisaldehyde and γ-himachalene.) Anethole is responsible for anise's characteristic odor and flavor. Other uses Builders of steam locomotives in Britain incorporated capsules of aniseed oil into white metal plain bearings so the distinctive smell would give warning in case of overheating. Anise can be made into a liquid scent and is used for both drag hunting and fishing. It is put on fishing lures to attract fish.	Biology and health sciences	Apiales	null
68686	https://en.wikipedia.org/wiki/Raman%20spectroscopy	Raman spectroscopy	Raman spectroscopy () (named after physicist C. V. Raman) is a spectroscopic technique typically used to determine vibrational modes of molecules, although rotational and other low-frequency modes of systems may also be observed. Raman spectroscopy is commonly used in chemistry to provide a structural fingerprint by which molecules can be identified. Raman spectroscopy relies upon inelastic scattering of photons, known as Raman scattering. A source of monochromatic light, usually from a laser in the visible, near infrared, or near ultraviolet range is used, although X-rays can also be used. The laser light interacts with molecular vibrations, phonons or other excitations in the system, resulting in the energy of the laser photons being shifted up or down. The shift in energy gives information about the vibrational modes in the system. Infrared spectroscopy typically yields similar yet complementary information. Typically, a sample is illuminated with a laser beam. Electromagnetic radiation from the illuminated spot is collected with a lens and sent through a monochromator. Elastic scattered radiation at the wavelength corresponding to the laser line (Rayleigh scattering) is filtered out by either a notch filter, edge pass filter, or a band pass filter, while the rest of the collected light is dispersed onto a detector. Spontaneous Raman scattering is typically very weak. As a result, for many years the main difficulty in collecting Raman spectra was separating the weak inelastically scattered light from the intense Rayleigh scattered laser light (referred to as "laser rejection"). Historically, Raman spectrometers used holographic gratings and multiple dispersion stages to achieve a high degree of laser rejection. In the past, photomultipliers were the detectors of choice for dispersive Raman setups, which resulted in long acquisition times. However, modern instrumentation almost universally employs notch or edge filters for laser rejection. Dispersive single-stage spectrographs (axial transmissive (AT) or Czerny–Turner (CT) monochromators) paired with CCD detectors are most common although Fourier transform (FT) spectrometers are also common for use with NIR lasers. The name "Raman spectroscopy" typically refers to vibrational Raman spectroscopy using laser wavelengths which are not absorbed by the sample. There are many other variations of Raman spectroscopy including surface-enhanced Raman, resonance Raman, tip-enhanced Raman, polarized Raman, stimulated Raman, transmission Raman, spatially-offset Raman, and hyper Raman. History Although the inelastic scattering of light was predicted by Adolf Smekal in 1923, it was not observed in practice until 1928. The Raman effect was named after one of its discoverers, the Indian scientist C. V. Raman, who observed the effect in organic liquids in 1928 together with K. S. Krishnan, and independently by Grigory Landsberg and Leonid Mandelstam in inorganic crystals. Raman won the Nobel Prize in Physics in 1930 for this discovery. The first observation of Raman spectra in gases was in 1929 by Franco Rasetti. Systematic pioneering theory of the Raman effect was developed by Czechoslovak physicist George Placzek between 1930 and 1934. The mercury arc became the principal light source, first with photographic detection and then with spectrophotometric detection. In the years following its discovery, Raman spectroscopy was used to provide the first catalog of molecular vibrational frequencies. Typically, the sample was held in a long tube and illuminated along its length with a beam of filtered monochromatic light generated by a gas discharge lamp. The photons that were scattered by the sample were collected through an optical flat at the end of the tube. To maximize the sensitivity, the sample was highly concentrated (1 M or more) and relatively large volumes (5 mL or more) were used. Theory The magnitude of the Raman effect correlates with the polarizability of the electrons in a molecule. It is a form of inelastic light scattering, where a photon excites the sample. This excitation puts the molecule into a virtual energy state for a short time before the photon is emitted. Inelastic scattering means that the energy of the emitted photon is of either lower or higher energy than the incident photon. After the scattering event, the sample is in a different rotational or vibrational state. For the total energy of the system to remain constant after the molecule moves to a new rovibronic (rotational–vibrational–electronic) state, the scattered photon shifts to a different energy, and therefore a different frequency. This energy difference is equal to that between the initial and final rovibronic states of the molecule. If the final state is higher in energy than the initial state, the scattered photon will be shifted to a lower frequency (lower energy) so that the total energy remains the same. This shift in frequency is called a Stokes shift, or downshift. If the final state is lower in energy, the scattered photon will be shifted to a higher frequency, which is called an anti-Stokes shift, or upshift. For a molecule to exhibit a Raman effect, there must be a change in its electric dipole-electric dipole polarizability with respect to the vibrational coordinate corresponding to the rovibronic state. The intensity of the Raman scattering is proportional to this polarizability change. Therefore, the Raman spectrum (scattering intensity as a function of the frequency shifts) depends on the rovibronic states of the molecule. The Raman effect is based on the interaction between the electron cloud of a sample and the external electric field of the monochromatic light, which can create an induced dipole moment within the molecule based on its polarizability. Because the laser light does not excite the molecule there can be no real transition between energy levels. The Raman effect should not be confused with emission (fluorescence or phosphorescence), where a molecule in an excited electronic state emits a photon and returns to the ground electronic state, in many cases to a vibrationally excited state on the ground electronic state potential energy surface. Raman scattering also contrasts with infrared (IR) absorption, where the energy of the absorbed photon matches the difference in energy between the initial and final rovibronic states. The dependence of Raman on the electric dipole-electric dipole polarizability derivative also differs from IR spectroscopy, which depends on the electric dipole moment derivative, the atomic polar tensor (APT). This contrasting feature allows rovibronic transitions that might not be active in IR to be analyzed using Raman spectroscopy, as exemplified by the rule of mutual exclusion in centrosymmetric molecules. Transitions which have large Raman intensities often have weak IR intensities and vice versa. If a bond is strongly polarized, a small change in its length such as that which occurs during a vibration has only a small resultant effect on polarization. Vibrations involving polar bonds (e.g. C-O , N-O , O-H) are therefore, comparatively weak Raman scatterers. Such polarized bonds, however, carry their electrical charges during the vibrational motion, (unless neutralized by symmetry factors), and this results in a larger net dipole moment change during the vibration, producing a strong IR absorption band. Conversely, relatively neutral bonds (e.g. C-C , C-H , C=C) suffer large changes in polarizability during a vibration. However, the dipole moment is not similarly affected such that while vibrations involving predominantly this type of bond are strong Raman scatterers, they are weak in the IR. A third vibrational spectroscopy technique, inelastic incoherent neutron scattering (IINS), can be used to determine the frequencies of vibrations in highly symmetric molecules that may be both IR and Raman inactive. The IINS selection rules, or allowed transitions, differ from those of IR and Raman, so the three techniques are complementary. They all give the same frequency for a given vibrational transition, but the relative intensities provide different information due to the different types of interaction between the molecule and the incoming particles, photons for IR and Raman, and neutrons for IINS. Raman shift Raman shifts are typically reported in wavenumbers, which have units of inverse length, as this value is directly related to energy. In order to convert between spectral wavelength and wavenumbers of shift in the Raman spectrum, the following formula can be used: where is the Raman shift expressed in wavenumber, is the excitation wavelength, and is the Raman spectrum wavelength. Most commonly, the unit chosen for expressing wavenumber in Raman spectra is inverse centimeters (cm−1). Since wavelength is often expressed in units of nanometers (nm), the formula above can scale for this unit conversion explicitly, giving Instrumentation Modern Raman spectroscopy nearly always involves the use of lasers as excitation light sources. Because lasers were not available until more than three decades after the discovery of the effect, Raman and Krishnan used a mercury lamp and photographic plates to record spectra. Early spectra took hours or even days to acquire due to weak light sources, poor sensitivity of the detectors and the weak Raman scattering cross-sections of most materials. Various colored filters and chemical solutions were used to select certain wavelength regions for excitation and detection but the photographic spectra were still dominated by a broad center line corresponding to Rayleigh scattering of the excitation source. Technological advances have made Raman spectroscopy much more sensitive, particularly since the 1980s. The most common modern detectors are now charge-coupled devices (CCDs). Photodiode arrays and photomultiplier tubes were common prior to the adoption of CCDs. The advent of reliable, stable, inexpensive lasers with narrow bandwidths has also had an impact. Lasers Raman spectroscopy requires a light source such as a laser. The resolution of the spectrum relies on the bandwidth of the laser source used. Generally shorter wavelength lasers give stronger Raman scattering due to the 4 increase in Raman scattering cross-sections, but issues with sample degradation or fluorescence may result. Continuous wave lasers are most common for normal Raman spectroscopy, but pulsed lasers may also be used. These often have wider bandwidths than their CW counterparts but are very useful for other forms of Raman spectroscopy such as transient, time-resolved and resonance Raman. Detectors Raman scattered light is typically collected and either dispersed by a spectrograph or used with an interferometer for detection by Fourier Transform (FT) methods. In many cases commercially available FT-IR spectrometers can be modified to become FT-Raman spectrometers. Detectors for dispersive Raman In most cases, modern Raman spectrometers use array detectors such as CCDs. Various types of CCDs exist which are optimized for different wavelength ranges. Intensified CCDs can be used for very weak signals and/or pulsed lasers. The spectral range depends on the size of the CCD and the focal length of spectrograph used. It was once common to use monochromators coupled to photomultiplier tubes. In this case the monochromator would need to be moved in order to scan through a spectral range. Detectors for FT–Raman FT–Raman is almost always used with NIR lasers and appropriate detectors must be used depending on the exciting wavelength. Germanium or Indium gallium arsenide (InGaAs) detectors are commonly used. Filters It is usually necessary to separate the Raman scattered light from the Rayleigh signal and reflected laser signal in order to collect high quality Raman spectra using a laser rejection filter. Notch or long-pass optical filters are typically used for this purpose. Before the advent of holographic filters it was common to use a triple-grating monochromator in subtractive mode to isolate the desired signal. This may still be used to record very small Raman shifts as holographic filters typically reflect some of the low frequency bands in addition to the unshifted laser light. However, Volume hologram filters are becoming more common which allow shifts as low as 5 cm−1 to be observed. Applications Raman spectroscopy is used in chemistry to identify molecules and study chemical bonding and intramolecular bonds. Because vibrational frequencies are specific to a molecule's chemical bonds and symmetry (the fingerprint region of organic molecules is in the wavenumber range 500–1,500 cm−1), Raman provides a fingerprint to identify molecules. For instance, Raman and IR spectra were used to determine the vibrational frequencies of SiO, Si2O2, and Si3O3 on the basis of normal coordinate analyses. Raman is also used to study the addition of a substrate to an enzyme. In solid-state physics, Raman spectroscopy is used to characterize materials, measure temperature, and find the crystallographic orientation of a sample. As with single molecules, a solid material can be identified by characteristic phonon modes. Information on the population of a phonon mode is given by the ratio of the Stokes and anti-Stokes intensity of the spontaneous Raman signal. Raman spectroscopy can also be used to observe other low frequency excitations of a solid, such as plasmons, magnons, and superconducting gap excitations. Distributed temperature sensing (DTS) uses the Raman-shifted backscatter from laser pulses to determine the temperature along optical fibers. The orientation of an anisotropic crystal can be found from the polarization of Raman-scattered light with respect to the crystal and the polarization of the laser light, if the crystal structure’s point group is known. In nanotechnology, a Raman microscope can be used to analyze nanowires to better understand their structures, and the radial breathing mode of carbon nanotubes is commonly used to evaluate their diameter. Raman active fibers, such as aramid and carbon, have vibrational modes that show a shift in Raman frequency with applied stress. Polypropylene fibers exhibit similar shifts. In solid state chemistry and the bio-pharmaceutical industry, Raman spectroscopy can be used to not only identify active pharmaceutical ingredients (APIs), but to identify their polymorphic forms, if more than one exist. For example, the drug Cayston (aztreonam), marketed by Gilead Sciences for cystic fibrosis, can be identified and characterized by IR and Raman spectroscopy. Using the correct polymorphic form in bio-pharmaceutical formulations is critical, since different forms have different physical properties, like solubility and melting point. Raman spectroscopy has a wide variety of applications in biology and medicine. It has helped confirm the existence of low-frequency phonons in proteins and DNA, promoting studies of low-frequency collective motion in proteins and DNA and their biological functions. Raman reporter molecules with olefin or alkyne moieties are being developed for tissue imaging with SERS-labeled antibodies. Raman spectroscopy has also been used as a noninvasive technique for real-time, in situ biochemical characterization of wounds. Multivariate analysis of Raman spectra has enabled development of a quantitative measure for wound healing progress. Spatially offset Raman spectroscopy (SORS), which is less sensitive to surface layers than conventional Raman, can be used to discover counterfeit drugs without opening their packaging, and to non-invasively study biological tissue. A reason why Raman spectroscopy is useful in biological applications is because its results often do not face interference from water molecules, due to the fact that they have permanent dipole moments, and as a result, the Raman scattering cannot be picked up on. This is a large advantage, specifically in biological applications. Raman spectroscopy also has a wide usage for studying biominerals. Lastly, Raman gas analyzers have many practical applications, including real-time monitoring of anesthetic and respiratory gas mixtures during surgery. Raman spectroscopy has been used in several research projects as a means to detect explosives from a safe distance using laser beams. Raman Spectroscopy is being further developed so it could be used in the clinical setting. Raman4Clinic is a European organization that is working on incorporating Raman Spectroscopy techniques in the medical field. They are currently working on different projects, one of them being monitoring cancer using bodily fluids such as urine and blood samples which are easily accessible. This technique would be less stressful on the patients than constantly having to take biopsies which are not always risk free. In photovoltaics, Raman spectroscopy has gained more interest in the past few years demonstrating high efficacy in delivering important properties for such materials. This includes optoelectronic and physicochemical properties such as open circuit voltage, efficiency, and crystalline structure. This has been demonstrated with several photovoltaic technologies, including kesterite-based, CIGS devices, Monocrystalline silicon cells, and perovskites devices. Art and cultural heritage Raman spectroscopy is an efficient and non-destructive way to investigate works of art and cultural heritage artifacts, in part because it is a non-invasive process which can be applied in situ. It can be used to analyze the corrosion products on the surfaces of artifacts (statues, pottery, etc.), which can lend insight into the corrosive environments experienced by the artifacts. The resulting spectra can also be compared to the spectra of surfaces that are cleaned or intentionally corroded, which can aid in determining the authenticity of valuable historical artifacts. It is capable of identifying individual pigments in paintings and their degradation products, which can provide insight into the working method of an artist in addition to aiding in authentication of paintings. It also gives information about the original state of the painting in cases where the pigments have degraded with age. Beyond the identification of pigments, extensive Raman microspectroscopic imaging has been shown to provide access to a plethora of trace compounds in Early Medieval Egyptian blue, which enable to reconstruct the individual "biography" of a colourant, including information on the type and provenance of the raw materials, synthesis and application of the pigment, and the ageing of the paint layer. In addition to paintings and artifacts, Raman spectroscopy can be used to investigate the chemical composition of historical documents (such as the Book of Kells), which can provide insight about the social and economic conditions when they were created. It also offers a noninvasive way to determine the best method of preservation or conservation of such cultural heritage artifacts, by providing insight into the causes behind deterioration. The IRUG (Infrared and Raman Users Group) Spectral Database is a rigorously peer-reviewed online database of IR and Raman reference spectra for cultural heritage materials such as works of art, architecture, and archaeological artifacts. The database is open for the general public to peruse, and includes interactive spectra for over a hundred different types of pigments and paints. Microspectroscopy Raman spectroscopy offers several advantages for microscopic analysis. Since it is a light scattering technique, specimens do not need to be fixed or sectioned. Raman spectra can be collected from a very small volume (< 1 μm in diameter, < 10 μm in depth); these spectra allow the identification of species present in that volume. Water does not generally interfere with Raman spectral analysis. Thus, Raman spectroscopy is suitable for the microscopic examination of minerals, materials such as polymers and ceramics, cells, proteins and forensic trace evidence. A Raman microscope begins with a standard optical microscope, and adds an excitation laser, a monochromator or polychromator, and a sensitive detector (such as a charge-coupled device (CCD), or photomultiplier tube (PMT)). FT-Raman has also been used with microscopes, typically in combination with near-infrared (NIR) laser excitation. Ultraviolet microscopes and UV enhanced optics must be used when a UV laser source is used for Raman microspectroscopy. In direct imaging (also termed global imaging or wide-field illumination), the whole field of view is examined for light scattering integrated over a small range of wavenumbers (Raman shifts). For instance, a wavenumber characteristic for cholesterol could be used to record the distribution of cholesterol within a cell culture. This technique is being used for the characterization of large-scale devices, mapping of different compounds and dynamics study. It has already been used for the characterization of graphene layers, J-aggregated dyes inside carbon nanotubes and multiple other 2D materials such as MoS2 and WSe2. Since the excitation beam is dispersed over the whole field of view, those measurements can be done without damaging the sample. The most common approach is hyperspectral imaging or chemical imaging, in which thousands of Raman spectra are acquired from all over the field of view by, for example, raster scanning of a focused laser beam through a sample. The data can be used to generate images showing the location and amount of different components. Having the full spectroscopic information available in every measurement spot has the advantage that several components can be mapped at the same time, including chemically similar and even polymorphic forms, which cannot be distinguished by detecting only one single wavenumber. Furthermore, material properties such as stress and strain, crystal orientation, crystallinity and incorporation of foreign ions into crystal lattices (e.g., doping, solid solution series) can be determined from hyperspectral maps. Taking the cell culture example, a hyperspectral image could show the distribution of cholesterol, as well as proteins, nucleic acids, and fatty acids. Sophisticated signal- and image-processing techniques can be used to ignore the presence of water, culture media, buffers, and other interferences. Because a Raman microscope is a diffraction-limited system, its spatial resolution depends on the wavelength of light, the numerical aperture of the focusing element, and — in the case of confocal microscopy — on the diameter of the confocal aperture. When operated in the visible to near-infrared range, a Raman microscope can achieve lateral resolutions of approx. 1 μm down to 250 nm, depending on the wavelength and type of objective lens (e.g., air vs. water or oil immersion lenses). The depth resolution (if not limited by the optical penetration depth of the sample) can range from 1–6 μm with the smallest confocal pinhole aperture to tens of micrometers when operated without a confocal pinhole. Depending on the sample, the high laser power density due to microscopic focussing can have the benefit of enhanced photobleaching of molecules emitting interfering fluorescence. However, the laser wavelength and laser power have to be carefully selected for each type of sample to avoid its degradation. Applications of Raman imaging range from materials sciences to biological studies. For each type of sample, the measurement parameters have to be individually optimized. For that reason, modern Raman microscopes are often equipped with several lasers offering different wavelengths, a set of objective lenses, and neutral density filters for tuning of the laser power reaching the sample. Selection of the laser wavelength mainly depends on optical properties of the sample and on the aim of the investigation. For example, Raman microscopy of biological and medical specimens is often performed using red to near-infrared excitation (e.g., 785 nm, or 1,064 nm wavelength). Due to typically low absorbances of biological samples in this spectral range, the risk of damaging the specimen as well as autofluorescence emission are reduced, and high penetration depths into tissues can be achieved. However, the intensity of Raman scattering at long wavelengths is low (owing to the ω4 dependence of Raman scattering intensity), leading to long acquisition times. On the other hand, resonance Raman imaging of single-cell algae at 532 nm (green) can specifically probe the carotenoid distribution within a cell by a using low laser power of ~5 μW and only 100 ms acquisition time. Raman scattering, specifically tip-enhanced Raman spectroscopy, produces high resolution hyperspectral images of single molecules, atoms, and DNA. Polarization dependence of Raman scattering Raman scattering is polarization sensitive and can provide detailed information on symmetry of Raman active modes. While conventional Raman spectroscopy identifies chemical composition, polarization effects on Raman spectra can reveal information on the orientation of molecules in single crystals and anisotropic materials, e.g. strained plastic sheets, as well as the symmetry of vibrational modes. Polarization–dependent Raman spectroscopy uses (plane) polarized laser excitation from a polarizer. The Raman scattered light collected is passed through a second polarizer (called the analyzer) before entering the detector. The analyzer is oriented either parallel or perpendicular to the polarization of the laser. Spectra acquired with the analyzer set at both perpendicular and parallel to the excitation plane can be used to calculate the depolarization ratio. Typically a polarization scrambler is placed between the analyzer and detector also.It is convenient in polarized Raman spectroscopy to describe the propagation and polarization directions using Porto's notation, described by and named after Brazilian physicist Sergio Pereira da Silva Porto. For isotropic solutions, the Raman scattering from each mode either retains the polarization of the laser or becomes partly or fully depolarized. If the vibrational mode involved in the Raman scattering process is totally symmetric then the polarization of the Raman scattering will be the same as that of the incoming laser beam. In the case that the vibrational mode is not totally symmetric then the polarization will be lost (scrambled) partially or totally, which is referred to as depolarization. Hence polarized Raman spectroscopy can provide detailed information as to the symmetry labels of vibrational modes. In the solid state, polarized Raman spectroscopy can be useful in the study of oriented samples such as single crystals. The polarizability of a vibrational mode is not equal along and across the bond. Therefore the intensity of the Raman scattering will be different when the laser's polarization is along and orthogonal to a particular bond axis. This effect can provide information on the orientation of molecules with a single crystal or material. The spectral information arising from this analysis is often used to understand macro-molecular orientation in crystal lattices, liquid crystals or polymer samples. Characterization of the symmetry of a vibrational mode The polarization technique is useful in understanding the connections between molecular symmetry, Raman activity, and peaks in the corresponding Raman spectra. Polarized light in one direction only gives access to some Raman–active modes, but rotating the polarization gives access to other modes. Each mode is separated according to its symmetry. The symmetry of a vibrational mode is deduced from the depolarization ratio ρ, which is the ratio of the Raman scattering with polarization orthogonal to the incident laser and the Raman scattering with the same polarization as the incident laser: Here is the intensity of Raman scattering when the analyzer is rotated 90 degrees with respect to the incident light's polarization axis, and the intensity of Raman scattering when the analyzer is aligned with the polarization of the incident laser. When polarized light interacts with a molecule, it distorts the molecule which induces an equal and opposite effect in the plane-wave, causing it to be rotated by the difference between the orientation of the molecule and the angle of polarization of the light wave. If , then the vibrations at that frequency are depolarized; meaning they are not totally symmetric. Raman Excitation Profile Analysis Resonance Raman selection rules can be explained by the Kramers-Heisenberg-Dirac (KHD) equation using the Albrecht A and B terms, as demonstrated. The KHD expression is conveniently linked to the polarizability of the molecule within its frame of reference. The polarizability operator connecting the initial and final states expresses the transition polarizability as a matrix element, as a function of the incidence frequency ω0. The directions x, y, and z in the molecular frame are represented by the Cartesian tensor ρ and σ here. Analyzing Raman excitation patterns requires the use of this equation, which is a sum-over-states expression for polarizability. This series of profiles illustrates the connection between a Raman active vibration's excitation frequency and intensity. This method takes into account sums over Franck-Condon's active vibrational states and provides insight into electronic absorption and emission spectra. Nevertheless, the work highlights a flaw in the sum-over-states method, especially for large molecules like visible chromophores, which are commonly studied in Raman spectroscopy. The difficulty arises from the potentially infinite number of intermediary steps needed. While lowering the sum at higher vibrational states can help tiny molecules get over this issue, larger molecules find it more challenging when there are more terms in the sum, particularly in the condensed phase when individual eigenstates cannot be resolved spectrally. To overcome this, two substitute techniques that do not require adding eigenstates can be considered. Among these two methods are available: the transform method. and Heller's time-dependent approach. The goal of both approaches is to take into consideration the frequency-dependent Raman cross-section σR(ω0) of a particular normal mode. Variants At least 25 variations of Raman spectroscopy have been developed. The usual purpose is to enhance the sensitivity (e.g., Surface-enhanced Raman spectroscopy (SERS)), to improve the spatial resolution (Raman microscopy), or to acquire very specific information (resonance Raman). Spontaneous (or far-field) Raman spectroscopy Terms such as spontaneous Raman spectroscopy or normal Raman spectroscopy summarize Raman spectroscopy techniques based on Raman scattering by using normal far-field optics as described above. Variants of normal Raman spectroscopy exist with respect to excitation-detection geometries, combination with other techniques, use of special (polarizing) optics and specific choice of excitation wavelengths for resonance enhancement. Correlative Raman imaging – Raman microscopy can be combined with complementary imaging methods, such as atomic force microscopy (Raman-AFM) and scanning electron microscopy (Raman-SEM) to compare Raman distribution maps with (or overlay them onto) topographical or morphological images, and to correlate Raman spectra with complementary physical or chemical information (e.g., gained by SEM-EDX). Resonance Raman spectroscopy – The excitation wavelength is matched to an electronic transition of the molecule or crystal, so that vibrational modes associated with the excited electronic state are greatly enhanced. This is useful for studying large molecules such as polypeptides, which might show hundreds of bands in "conventional" Raman spectra. It is also useful for associating normal modes with their observed frequency shifts. Angle-resolved Raman spectroscopy – Not only are standard Raman results recorded but also the angle with respect to the incident laser. If the orientation of the sample is known then detailed information about the phonon dispersion relation can also be gleaned from a single test. Optical tweezers Raman spectroscopy (OTRS) – Used to study individual particles, and even biochemical processes in single cells trapped by optical tweezers. Spatially offset Raman spectroscopy (SORS) – The Raman scattering beneath an obscuring surface is retrieved from a scaled subtraction of two spectra taken at two spatially offset points. Raman optical activity (ROA) – Measures vibrational optical activity by means of a small difference in the intensity of Raman scattering from chiral molecules in right- and left-circularly polarized incident light or, equivalently, a small circularly polarized component in the scattered light. Transmission Raman – Allows probing of a significant bulk of a turbid material, such as powders, capsules, living tissue, etc. It was largely ignored following investigations in the late 1960s (Schrader and Bergmann, 1967) but was rediscovered in 2006 as a means of rapid assay of pharmaceutical dosage forms. There are medical diagnostic applications particularly in the detection of cancer. Micro-cavity substrates – A method that improves the detection limit of conventional Raman spectra using micro-Raman in a micro-cavity coated with reflective Au or Ag. The micro-cavity has a radius of several micrometers and enhances the entire Raman signal by providing multiple excitations of the sample and couples the forward-scattered Raman photons toward the collection optics in the back-scattered Raman geometry. Stand-off remote Raman – In standoff Raman, the sample is measured at a distance from the Raman spectrometer, usually by using a telescope for light collection. Remote Raman spectroscopy was proposed in the 1960s and initially developed for the measurement of atmospheric gases. The technique was extended In 1992 by Angel et al. for standoff Raman detection of hazardous inorganic and organic compounds. X-ray Raman scattering – Measures electronic transitions rather than vibrations. Enhanced (or near-field) Raman spectroscopy Enhancement of Raman scattering is achieved by local electric-field enhancement by optical near-field effects (e.g. localized surface plasmons). Surface-enhanced Raman spectroscopy (SERS) – Normally done in a silver or gold colloid or a substrate containing silver or gold. Surface plasmons of silver and gold are excited by the laser, resulting in an increase in the electric fields surrounding the metal. Given that Raman intensities are proportional to the electric field, there is large increase in the measured signal (by up to 1011). This effect was originally observed by Martin Fleischmann but the prevailing explanation was proposed by Van Duyne in 1977. A comprehensive theory of the effect was given by Lombardi and Birke. Surface-enhanced resonance Raman spectroscopy (SERRS) – A combination of SERS and resonance Raman spectroscopy that uses proximity to a surface to increase Raman intensity, and excitation wavelength matched to the maximum absorbance of the molecule being analysed. Tip-enhanced Raman spectroscopy (TERS) – TERS combines the chemical sensitivity of SERS with the high spatial resolution of scanning probe microscopy techniques, enabling chemical imaging of surfaces at the nanometre length-scale with high detection sensitivity. It uses a metallic (usually silver-/gold-coated AFM or STM) tip to enhance the Raman signals of molecules situated in its vicinity. The spatial resolution is approximately the size of the tip apex (20–30 nm). TERS has been shown to have sensitivity down to the single molecule level and holds some promise for bioanalysis applications and DNA sequencing. TERS was used to image the vibrational normal modes of single molecules. Surface plasmon polariton enhanced Raman scattering (SPPERS) – This approach exploits apertureless metallic conical tips for near field excitation of molecules. This technique differs from the TERS approach due to its inherent capability of suppressing the background field. In fact, when an appropriate laser source impinges on the base of the cone, a TM0 mode (polaritonic mode) can be locally created, namely far away from the excitation spot (apex of the tip). The mode can propagate along the tip without producing any radiation field up to the tip apex where it interacts with the molecule. In this way, the focal plane is separated from the excitation plane by a distance given by the tip length, and no background plays any role in the Raman excitation of the molecule. Non-linear Raman spectroscopy Raman signal enhancements are achieved through non-linear optical effects, typically realized by mixing two or more wavelengths emitted by spatially and temporally synchronized pulsed lasers. Hyper Raman – A non-linear effect in which the vibrational modes interact with the second harmonic of the excitation beam. This requires very high power, but allows the observation of vibrational modes that are normally "silent". It frequently relies on SERS-type enhancement to boost the sensitivity. Stimulated Raman spectroscopy (SRS) – A pump-probe technique, where a spatially coincident, two color pulse (with polarization either parallel or perpendicular) transfers the population from ground to a rovibrationally excited state. If the difference in energy corresponds to an allowed Raman transition, scattered light will correspond to loss or gain in the pump beam. Inverse Raman spectroscopy – A synonym for stimulated Raman loss spectroscopy. Coherent anti-Stokes Raman spectroscopy (CARS) – Two laser beams are used to generate a coherent anti-Stokes frequency beam, which can be enhanced by resonance. Morphologically-Directed Raman spectroscopy Morphologically Directed Raman Spectroscopy (MDRS) combines automated particle imaging and Raman microspectroscopy into a singular integrated platform in order to provide particle size, shape, and chemical identification. Automated particle imaging determines the particle size and shape distributions of components within a blended sample from images of individual particles. The information gathered from automated particle imaging is then utilized to direct the Raman spectroscopic analysis. The Raman spectroscopic analytical process is performed on a randomly-selected subset of the particles, allowing chemical identification of the sample’s multiple components. Tens of thousands of particles can be imaged in a matter of minutes using the MDRS method, making the process ideal for forensic analysis and investigating counterfeit pharmaceuticals and subsequent adjudications.	Physical sciences	Spectroscopy	Chemistry
68756	https://en.wikipedia.org/wiki/Human%20tooth	Human tooth	Human teeth function to mechanically break down items of food by cutting and crushing them in preparation for swallowing and digesting. As such, they are considered part of the human digestive system. Humans have four types of teeth: incisors, canines, premolars, and molars, which each have a specific function. The incisors cut the food, the canines tear the food and the molars and premolars crush the food. The roots of teeth are embedded in the maxilla (upper jaw) or the mandible (lower jaw) and are covered by gums. Teeth are made of multiple tissues of varying density and hardness. Humans, like most other mammals, are diphyodont, meaning that they develop two sets of teeth. The first set, deciduous teeth, also called "primary teeth", "baby teeth", or "milk teeth", normally eventually contains 20 teeth. Primary teeth typically start to appear ("erupt") around six months of age and this may be distracting and/or painful for the infant. However, some babies are born with one or more visible teeth, known as neonatal teeth or "natal teeth". Structure Dental anatomy is dedicated to the study of tooth structure. The development, appearance, and classification of teeth fall within its field of study, though dental occlusion, or contact between teeth, does not. Dental anatomy is also a taxonomic science as it is concerned with the naming of teeth and their structures. This information serves a practical purpose for dentists, enabling them to easily identify and describe teeth and structures during treatment. The anatomic crown of a tooth is the area covered in enamel above the cementoenamel junction (CEJ) or "neck" of the tooth. Most of the crown is composed of dentin ("dentine" in British English) with the pulp chamber inside. The crown is within bone before eruption. After eruption, it is almost always visible. The anatomic root is found below the CEJ and is covered with cementum. As with the crown, dentin composes most of the root, which normally has pulp canals. Canines and most premolars, except for maxillary first premolars, usually have one root. Maxillary first premolars and mandibular molars usually have two roots. Maxillary molars usually have three roots. Additional roots are referred to as supernumerary roots. Humans usually have 20 primary (deciduous, "baby" or "milk") teeth and 32 permanent (adult) teeth. Teeth are classified as incisors, canines, premolars (also called bicuspids), and molars. Incisors are primarily used for cutting, canines are for tearing, and molars serve for grinding. Most teeth have identifiable features that distinguish them from others. There are several different notation systems to refer to a specific tooth. The three most common systems are the FDI World Dental Federation notation (ISO 3950), the Universal Numbering System, and the Palmer notation. The FDI system is used worldwide, the Universal only in the United States, while the older Palmer notation still has some adherents only in the United Kingdom. Primary teeth Among deciduous (primary) teeth, ten are found in the maxilla (upper jaw) and ten in the mandible (lower jaw), for a total of 20. The dental formula for primary teeth in humans is . In the primary set of teeth, in addition to the canines there are two types of incisors—centrals and laterals—and two types of molars—first and second. All primary teeth are normally later replaced with their permanent counterparts. Permanent teeth Among permanent teeth, 16 are found in the maxilla and 16 in the mandible, for a total of 32. The dental formula is . Permanent human teeth are numbered in a boustrophedonic sequence. The maxillary teeth are the maxillary central incisors (teeth 8 and 9 in the diagram), maxillary lateral incisors (7 and 10), maxillary canines (6 and 11), maxillary first premolars (5 and 12), maxillary second premolars (4 and 13), maxillary first molars (3 and 14), maxillary second molars (2 and 15), and maxillary third molars (1 and 16). The mandibular teeth are the mandibular central incisors (24 and 25), mandibular lateral incisors (23 and 26), mandibular canines (22 and 27), mandibular first premolars (21 and 28), mandibular second premolars (20 and 29), mandibular first molars (19 and 30), mandibular second molars (18 and 31), and mandibular third molars (17 and 32). Third molars are commonly called "wisdom teeth" and usually emerge at ages 17 to 25. These molars may never erupt into the mouth or form at all. When they do form, they often must be removed. If any additional teeth form—for example, fourth and fifth molars, which are rare—they are referred to as supernumerary teeth (hyperdontia). Development of fewer than the usual number of teeth is called hypodontia. There are small differences between the teeth of males and females, with male teeth along with the male jaw tending to be larger on average than female teeth and jaw. There are also differences in the internal dental tissue proportions, with male teeth consisting of proportionately more dentine while female teeth have proportionately more enamel. Parts Enamel Enamel is the hardest and most highly mineralized substance of the body. It has its origin from oral ectoderm. It is one of the four major tissues which make up the tooth, along with dentin, cementum, and dental pulp. It is normally visible and must be supported by underlying dentin. 96% of enamel consists of mineral, with water and organic material comprising the rest. The normal color of enamel varies from light yellow to grayish white. At the edges of teeth where there is no dentin underlying the enamel, the color sometimes has a slightly blue tone. Since enamel is semitranslucent, the color of dentin and any restorative dental material underneath the enamel strongly affects the appearance of a tooth. Enamel varies in thickness over the surface of the tooth and is often thickest at the cusp, up to 2.5mm, and thinnest at its border, which is seen clinically as the CEJ. The wear rate of enamel, called attrition, is 8 micrometers a year from normal factors. Enamel's primary mineral is hydroxyapatite, which is a crystalline calcium phosphate. The large amount of minerals in enamel accounts not only for its strength but also for its brittleness. Dentin, which is less mineralized and less brittle, compensates for enamel and is necessary as a support. Unlike dentin and bone, enamel does not contain collagen. Proteins of note in the development of enamel are ameloblastins, amelogenins, enamelins and tuftelins. It is believed that they aid in the development of enamel by serving as framework support, among other functions. In rare circumstances enamel can fail to form, leaving the underlying dentine exposed on the surface. Dentin Dentin is the substance between enamel or cementum and the pulp chamber. It is secreted by the odontoblasts of the dental pulp. The formation of dentin is known as dentinogenesis. The porous, yellow-hued material is made up of 70% inorganic materials, 20% organic materials, and 10% water by weight. Because it is softer than enamel, it decays more rapidly and is subject to severe cavities if not properly treated, but dentin still acts as a protective layer and supports the crown of the tooth. Dentin is a mineralized connective tissue with an organic matrix of collagenous proteins. Dentin has microscopic channels, called dentinal tubules, which radiate outward through the dentin from the pulp cavity to the exterior cementum or enamel border. The diameter of these tubules range from 2.5 μm near the pulp, to 1.2 μm in the midportion, and 900 nm near the dentino-enamel junction. Although they may have tiny side-branches, the tubules do not intersect with each other. Their length is dictated by the radius of the tooth. The three dimensional configuration of the dentinal tubules is genetically determined. There are three types of dentin, primary, secondary and tertiary. Secondary dentin is a layer of dentin produced after root formation and continues to form with age. Tertiary dentin is created in response to stimulus, such as cavities and tooth wear. Cementum Cementum is a specialized bone like substance covering the root of a tooth. It is approximately 45% inorganic material (mainly hydroxyapatite), 33% organic material (mainly collagen) and 22% water. Cementum is excreted by cementoblasts within the root of the tooth and is thickest at the root apex. Its coloration is yellowish and it is softer than dentin and enamel. The principal role of cementum is to serve as a medium by which the periodontal ligaments can attach to the tooth for stability. At the cement to enamel junction, the cementum is acellular due to its lack of cellular components, and this acellular type covers at least ⅔ of the root. The more permeable form of cementum, cellular cementum, covers about ⅓ of the root apex. Dental pulp The dental pulp is the central part of the tooth filled with soft connective tissue. This tissue contains blood vessels and nerves that enter the tooth from a hole at the apex of the root. Along the border between the dentin and the pulp are odontoblasts, which initiate the formation of dentin. Other cells in the pulp include fibroblasts, preodontoblasts, macrophages and T lymphocytes. The pulp is commonly called "the nerve" of the tooth. Development Tooth development is the complex process by which teeth form from embryonic cells, grow, and erupt into the mouth. Although many diverse species have teeth, their development is largely the same as in humans. For human teeth to have a healthy oral environment, enamel, dentin, cementum, and the periodontium must all develop during appropriate stages of fetal development. Primary teeth start to form in the development of the embryo between the sixth and eighth weeks, and permanent teeth begin to form in the twentieth week. If teeth do not start to develop at or near these times, they will not develop at all. A significant amount of research has focused on determining the processes that initiate tooth development. It is widely accepted that there is a factor within the tissues of the first pharyngeal arch that is necessary for the development of teeth. Tooth development is commonly divided into the following stages: the bud stage, the cap, the bell, and finally maturation. The staging of tooth development is an attempt to categorize changes that take place along a continuum; frequently it is difficult to decide what stage should be assigned to a particular developing tooth. This determination is further complicated by the varying appearance of different histologic sections of the same developing tooth, which can appear to be different stages. The tooth bud (sometimes called the tooth germ) is an aggregation of cells that eventually forms a tooth. It is organized into three parts: the enamel organ, the dental papilla and the dental follicle. The enamel organ is composed of the outer enamel epithelium, inner enamel epithelium, stellate reticulum and stratum intermedium. These cells give rise to ameloblasts, which produce enamel and the reduced enamel epithelium. The growth of cervical loop cells into the deeper tissues forms Hertwig's Epithelial Root Sheath, which determines a tooth's root shape. The dental papilla contains cells that develop into odontoblasts, which are dentin-forming cells. Additionally, the junction between the dental papilla and inner enamel epithelium determines the crown shape of a tooth. The dental follicle gives rise to three important cells: cementoblasts, osteoblasts, and fibroblasts. Cementoblasts form the cementum of a tooth. Osteoblasts give rise to the alveolar bone around the roots of teeth. Fibroblasts develop the periodontal ligaments which connect teeth to the alveolar bone through cementum. Eruption Tooth eruption in humans is a process in tooth development in which the teeth enter the mouth and become visible. Current research indicates that the periodontal ligaments play an important role in tooth eruption. Primary teeth erupt into the mouth from around six months until two years of age. These teeth are the only ones in the mouth until a person is about six years old. At that time, the first permanent tooth erupts. This stage, during which a person has a combination of primary and permanent teeth, is known as the mixed stage. The mixed stage lasts until the last primary tooth is lost and the remaining permanent teeth erupt into the mouth. There have been many theories about the cause of tooth eruption. One theory proposes that the developing root of a tooth pushes it into the mouth. Another, known as the cushioned hammock theory, resulted from microscopic study of teeth, which was thought to show a ligament around the root. It was later discovered that the "ligament" was merely an artifact created in the process of preparing the slide. Currently, the most widely held belief is that the periodontal ligaments provide the main impetus for the process. The onset of primary tooth loss has been found to correlate strongly with somatic and psychological criteria of school readiness. Supporting structures The periodontium is the supporting structure of a tooth, helping to attach the tooth to surrounding tissues and to allow sensations of touch and pressure. It consists of the cementum, periodontal ligaments, alveolar bone, and gingiva. Of these, cementum is the only one that is a part of a tooth. Periodontal ligaments connect the alveolar bone to the cementum. Alveolar bone surrounds the roots of teeth to provide support and creates what is commonly called an alveolus, or "socket". Lying over the bone is the gingiva or gum, which is readily visible in the mouth. Periodontal ligaments The periodontal ligament is a specialized connective tissue that attaches the cementum of a tooth to the alveolar bone. This tissue covers the root of the tooth within the bone. Each ligament has a width of 0.15–0.38mm, but this size decreases over time. The functions of the periodontal ligaments include attachment of the tooth to the bone, support for the tooth, formation and resorption of bone during tooth movement, sensation, and eruption. The cells of the periodontal ligaments include osteoblasts, osteoclasts, fibroblasts, macrophages, cementoblasts, and epithelial cell rests of Malassez. Consisting of mostly Type I and III collagen, the fibers are grouped in bundles and named according to their location. The groups of fibers are named alveolar crest, horizontal, oblique, periapical, and interradicular fibers. The nerve supply generally enters from the bone apical to the tooth and forms a network around the tooth toward the crest of the gingiva. When pressure is exerted on a tooth, such as during chewing or biting, the tooth moves slightly in its socket and puts tension on the periodontal ligaments. The nerve fibers can then send the information to the central nervous system for interpretation. Alveolar bone The alveolar bone is the bone of the jaw which forms the alveolus around teeth. Like any other bone in the human body, alveolar bone is modified throughout life. Osteoblasts create bone and osteoclasts destroy it, especially if force is placed on a tooth. As is the case when movement of teeth is attempted through orthodontics, an area of bone under compressive force from a tooth moving toward it has a high osteoclast level, resulting in bone resorption. An area of bone receiving tension from periodontal ligaments attached to a tooth moving away from it has a high number of osteoblasts, resulting in bone formation. Gingiva The gingiva ("gums") is the mucosal tissue that overlays the jaws. There are three different types of epithelium associated with the gingiva: gingival, junctional, and sulcular epithelium. These three types form from a mass of epithelial cells known as the epithelial cuff between the tooth and the mouth. The gingival epithelium is not associated directly with tooth attachment and is visible in the mouth. The junctional epithelium, composed of the basal lamina and hemidesmosomes, forms an attachment to the tooth. The sulcular epithelium is nonkeratinized stratified squamous tissue on the gingiva which touches but is not attached to the tooth. Tooth decay Plaque Plaque is a biofilm consisting of large quantities of various bacteria that form on teeth. If not removed regularly, plaque buildup can lead to periodontal problems such as gingivitis. Given time, plaque can mineralize along the gingiva, forming tartar. The microorganisms that form the biofilm are almost entirely bacteria (mainly streptococcus and anaerobes), with the composition varying by location in the mouth. Streptococcus mutans is the most important bacterium associated with dental caries. Certain bacteria in the mouth live off the remains of foods, especially sugars and starches. In the absence of oxygen they produce lactic acid, which dissolves the calcium and phosphorus in the enamel. This process, known as "demineralisation", leads to tooth destruction. Saliva gradually neutralises the acids, which causes the pH of the tooth surface to rise above the critical pH, typically considered to be 5.5. This causes remineralisation, the return of the dissolved minerals to the enamel. If there is sufficient time between the intake of foods then the impact is limited and the teeth can repair themselves. Saliva is unable to penetrate through plaque, however, to neutralize the acid produced by the bacteria. Caries (cavities) Dental caries (cavities), described as "tooth decay", is an infectious disease which damages the structures of teeth. The disease can lead to pain, tooth loss, and infection. Dental caries has a long history, with evidence showing the disease was present in the Bronze, Iron, and Middle ages but also prior to the neolithic period. The largest increases in the prevalence of caries have been associated with diet changes. Today, caries remains one of the most common diseases throughout the world. In the United States, dental caries is the most common chronic childhood disease, being at least five times more common than asthma. Countries that have experienced an overall decrease in cases of tooth decay continue to have a disparity in the distribution of the disease. Among children in the United States and Europe, 60–80% of cases of dental caries occur in 20% of the population. Tooth decay is caused by certain types of acid-producing bacteria which cause the most damage in the presence of fermentable carbohydrates such as sucrose, fructose, and glucose. The resulting acidic levels in the mouth affect teeth because a tooth's special mineral content causes it to be sensitive to low pH. Depending on the extent of tooth destruction, various treatments can be used to restore teeth to proper form, function, and aesthetics, but there is no known method to regenerate large amounts of tooth structure. Instead, dental health organizations advocate preventive and prophylactic measures, such as regular oral hygiene and dietary modifications, to avoid dental caries. Tooth care Oral hygiene Oral hygiene is the practice of keeping the mouth clean and is a means of preventing dental caries, gingivitis, periodontal disease, bad breath, and other dental disorders. It consists of both professional and personal care. Regular cleanings, usually done by dentists and dental hygienists, remove tartar (mineralized plaque) that may develop even with careful brushing and flossing. Professional cleaning includes tooth scaling, using various instruments or devices to loosen and remove deposits from teeth. The purpose of cleaning teeth is to remove plaque, which consists mostly of bacteria. Healthcare professionals recommend regular brushing twice a day (in the morning and in the evening, or after meals) in order to prevent formation of plaque and tartar. A toothbrush is able to remove most plaque, except in areas between teeth. As a result, flossing is also considered a necessity to maintain oral hygiene. When used correctly, dental floss removes plaque from between teeth and at the gum line, where periodontal disease often begins and could develop caries. Electric toothbrushes are a popular aid to oral hygiene. A user without disabilities, with proper training in manual brushing, and with good motivation, can achieve standards of oral hygiene at least as satisfactory as the best electric brushes, but untrained users rarely achieve anything of the kind. Not all electric toothbrushes are equally effective and even a good design needs to be used properly for best effect, but: "Electric toothbrushes tend to help people who are not as good at cleaning teeth and as a result have had oral hygiene problems." The most important advantage of electric toothbrushes is their ability to aid people with dexterity difficulties, such as those associated with rheumatoid arthritis. Protective treatments Fluoride therapy is often recommended to protect against dental caries. Water fluoridation and fluoride supplements decrease the incidence of dental caries. Fluoride helps prevent dental decay by binding to the hydroxyapatite crystals in enamel. The incorporated fluoride makes enamel more resistant to demineralization and thus more resistant to decay. Topical fluoride, such as a fluoride toothpaste or mouthwash, is also recommended to protect teeth surfaces. Many dentists include application of topical fluoride solutions as part of routine cleanings. Dental sealants are another preventive therapy often used to provide a barrier to bacteria and decay on the surface of teeth. Sealants can last up to ten years and are primarily used on the biting surfaces of molars of children and young adults, especially those who may have difficulty brushing and flossing effectively. Sealants are applied in a dentist's office, sometimes by a dental hygienist, in a procedure similar in technique and cost to a fluoride application. Restorations After a tooth has been damaged or destroyed, restoration of the missing structure can be achieved with a variety of treatments. Restorations may be created from a variety of materials, including glass ionomer, amalgam, gold, porcelain, and composite. Small restorations placed inside a tooth are referred to as "intracoronal restorations". These restorations may be formed directly in the mouth or may be cast using the lost-wax technique, such as for some inlays and onlays. When larger portions of a tooth are lost, an "extracoronal restoration" may be fabricated, such as an artificial crown or a veneer, to restore the involved tooth. When a tooth is lost, dentures, bridges, or implants may be used as replacements. Dentures are usually the least costly whereas implants are usually the most expensive. Dentures may replace complete arches of the mouth or only a partial number of teeth. Bridges replace smaller spaces of missing teeth and use adjacent teeth to support the restoration. Dental implants may be used to replace a single tooth or a series of teeth. Though implants are the most expensive treatment option, they are often the most desirable restoration because of their aesthetics and function. To improve the function of dentures, implants may be used as support. Abnormalities Tooth abnormalities may be categorized according to whether they have environmental or developmental causes. While environmental abnormalities may appear to have an obvious cause, there may not appear to be any known cause for some developmental abnormalities. Environmental forces may affect teeth during development, destroy tooth structure after development, discolor teeth at any stage of development, or alter the course of tooth eruption. Developmental abnormalities most commonly affect the number, size, shape, and structure of teeth. Environmental Alteration during tooth development Tooth abnormalities caused by environmental factors during tooth development have long-lasting effects. Enamel and dentin do not regenerate after they mineralize initially. Enamel hypoplasia is a condition in which the amount of enamel formed is inadequate. This results either in pits and grooves in areas of the tooth or in widespread absence of enamel. Diffuse opacities of enamel does not affect the amount of enamel but changes its appearance. Affected enamel has a different translucency than the rest of the tooth. Demarcated opacities of enamel have sharp boundaries where the translucency decreases and manifest a white, cream, yellow, or brown color. All these may be caused by nutritional factors, an exanthematous disease (chicken pox, congenital syphilis), undiagnosed and untreated celiac disease, hypocalcemia, dental fluorosis, birth injury, preterm birth, infection or trauma from a deciduous tooth. Dental fluorosis is a condition which results from ingesting excessive amounts of fluoride and leads to teeth which are spotted, yellow, brown, black or sometimes pitted. In most cases, the enamel defects caused by celiac disease, which may be the only manifestation of this disease in the absence of any other symptoms or signs, are not recognized and mistakenly attributed to other causes, such as fluorosis. Enamel hypoplasia resulting from syphilis is frequently referred to as Hutchinson's teeth, which is considered one part of Hutchinson's triad. Turner's hypoplasia is a portion of missing or diminished enamel on a permanent tooth usually from a prior infection of a nearby primary tooth. Hypoplasia may also result from antineoplastic therapy. Destruction after development Tooth destruction from processes other than dental caries is considered a normal physiologic process but may become severe enough to become a pathologic condition. Attrition is the loss of tooth structure by mechanical forces from opposing teeth. Attrition initially affects the enamel and, if unchecked, may proceed to the underlying dentin. Abrasion is the loss of tooth structure by mechanical forces from a foreign element. If this force begins at the cementoenamel junction, then progression of tooth loss can be rapid since enamel is very thin in this region of the tooth. A common source of this type of tooth wear is excessive force when using a toothbrush. Erosion is the loss of tooth structure due to chemical dissolution by acids not of bacterial origin. Signs of tooth destruction from erosion is a common characteristic in the mouths of people with bulimia since vomiting results in exposure of the teeth to gastric acids. Another important source of erosive acids are from frequent sucking of lemon juice. Abfraction is the loss of tooth structure from flexural forces. As teeth flex under pressure, the arrangement of teeth touching each other, known as occlusion, causes tension on one side of the tooth and compression on the other side of the tooth. This is believed to cause V-shaped depressions on the side under tension and C-shaped depressions on the side under compression. When tooth destruction occurs at the roots of teeth, the process is referred to as internal resorption, when caused by cells within the pulp, or external resorption, when caused by cells in the periodontal ligament. Discoloration Discoloration of teeth may result from bacteria stains, tobacco, tea, coffee, foods with an abundance of chlorophyll, restorative materials, and medications. Stains from bacteria may cause colors varying from green to black to orange. Green stains also result from foods with chlorophyll or excessive exposure to copper or nickel. Amalgam, a common dental restorative material, may turn adjacent areas of teeth black or gray. Long term use of chlorhexidine, a mouthwash, may encourage extrinsic stain formation near the gingiva on teeth. This is usually easy for a hygienist to remove. Systemic disorders also can cause tooth discoloration. Congenital erythropoietic porphyria causes porphyrins to be deposited in teeth, causing a red-brown coloration. Blue discoloration may occur with alkaptonuria and rarely with Parkinson's disease. Erythroblastosis fetalis and biliary atresia are diseases which may cause teeth to appear green from the deposition of biliverdin. Also, trauma may change a tooth to a pink, yellow, or dark gray color. Pink and red discolorations are also associated in patients with lepromatous leprosy. Some medications, such as tetracycline antibiotics, may become incorporated into the structure of a tooth, causing intrinsic staining of the teeth. Alteration of eruption Tooth eruption may be altered by some environmental factors. When eruption is prematurely stopped, the tooth is said to be impacted. The most common cause of tooth impaction is lack of space in the mouth for the tooth. Other causes may be tumors, cysts, trauma, and thickened bone or soft tissue. Tooth ankylosis occurs when the tooth has already erupted into the mouth but the cementum or dentin has fused with the alveolar bone. This may cause a person to retain their primary tooth instead of having it replaced by a permanent one. A technique for altering the natural progression of eruption is employed by orthodontists who wish to delay or speed up the eruption of certain teeth for reasons of space maintenance or otherwise preventing crowding and/or spacing. If a primary tooth is extracted before its succeeding permanent tooth's root reaches ⅓ of its total growth, the eruption of the permanent tooth will be delayed. Conversely, if the roots of the permanent tooth are more than ⅔ complete, the eruption of the permanent tooth will be accelerated. Between ⅓ and ⅔, it is unknown exactly what will occur to the speed of eruption. Developmental Abnormality in number Anodontia is the total lack of tooth development. Hyperdontia is the presence of a higher-than-normal number of teeth. Hypodontia is the lack of development of one or more teeth. Oligodontia may be used to describe the absence of 6 or more teeth. Some systemic disorders which may result in hyperdontia include Apert syndrome, cleidocranial dysostosis, Crouzon syndrome, Ehlers–Danlos syndrome, Gardner's syndrome, and Sturge–Weber syndrome. Some systemic disorders which may result in hypodontia include Crouzon syndrome, Ectodermal dysplasia, Ehlers–Danlos syndrome, and Gorlin syndrome. Abnormality in size Microdontia is a condition where teeth are smaller than the usual size. Macrodontia is where teeth are larger than the usual size. Microdontia of a single tooth is more likely to occur in a maxillary lateral incisor. The second most likely tooth to have microdontia are third molars. Macrodontia of all the teeth is known to occur in pituitary gigantism and pineal hyperplasia. It may also occur on one side of the face in cases of hemifacial hyperplasia. Abnormality in shape Gemination occurs when a developing tooth incompletely splits into the formation of two teeth. Fusion is the union of two adjacent teeth during development. Concrescence is the fusion of two separate teeth only in their cementum. Accessory cusps are additional cusps on a tooth and may manifest as a Talon cusp, Cusp of Carabelli, or Dens evaginatus. Dens invaginatus, also called Dens in dente, is a deep invagination in a tooth causing the appearance of a tooth within a tooth. Ectopic enamel is enamel found in an unusual location, such as the root of a tooth. Taurodontism is a condition where the body of the tooth and pulp chamber is enlarged, and is associated with Klinefelter syndrome, Tricho-dento-osseous syndrome, Triple X syndrome, and XYY syndrome. Hypercementosis is excessive formation of cementum, which may result from trauma, inflammation, acromegaly, rheumatic fever, and Paget's disease of bone. A dilaceration is a bend in the root which may have been caused by trauma to the tooth during formation. Supernumerary roots is the presence of a greater number of roots on a tooth than expected Cleft lip and palate and their association with dental anomalies There are many types of dental anomalies seen in cleft lip and palate (CLP) patients. Both sets of dentition may be affected; however, they are commonly seen in the affected side. Most frequently, missing teeth, supernumerary or discoloured teeth can be seen; however, enamel dysplasia, discolouration and delayed root development are also common. In children with cleft lip and palate, the lateral incisor in the alveolar cleft region has the highest prevalence of dental developmental disorders; this condition may be a cause of tooth crowding. This is important to consider in order to correctly plan treatment keeping in mind considerations for function and aesthetics. By correctly coordinating management invasive treatment procedures can be prevented resulting in successful and conservative treatment. There have been a plethora of research studies to calculate prevalence of certain dental anomalies in CLP populations however a variety of results have been obtained. In a study evaluating dental anomalies in Brazilian cleft patients, male patients had a higher incidence of CLP, agenesis, and supernumerary teeth than did female patients. In cases of complete CLP, the left maxillary lateral incisor was the most commonly absent tooth. Supernumerary teeth were typically located distal to the cleft. In a study of Jordanian subjects, the prevalence of dental anomaly was higher in CLP patients than in normal subjects. Missing teeth were observed in 66.7% of patients, with maxillary lateral incisor as the most frequently affected tooth. Supernumerary teeth were observed in 16.7% of patients; other findings included microdontia (37%), taurodontism (70.5%), transposition or ectopic teeth (30.8%), dilacerations (19.2%), and hypoplasia (30.8%). The incidence of microdontia, dilaceration, and hypoplasia was significantly higher in bilateral CLP patients than in unilateral CLP patients, and none of the anomalies showed any significant sexual dimorphism. It is therefore evident that patients with cleft lip and palate may present with a variety of dental anomalies. It is essential to assess the patient both clinically and radiographically in order to correctly treat and prevent progression of any dental problems. It is also useful to note that patients with a cleft lip and palate automatically score a 5 on the IOTN ( index for orthodontic need) and therefore are eligible for orthodontic treatment, liaising with an orthodontist is vital in order coordinate and plan treatment successfully. Abnormality in structure Amelogenesis imperfecta is a condition in which enamel does not form properly or at all. Dentinogenesis imperfecta is a condition in which dentin does not form properly and is sometimes associated with osteogenesis imperfecta. Dentin dysplasia is a disorder in which the roots and pulp of teeth may be affected. Regional odontodysplasia is a disorder affecting enamel, dentin, and pulp and causes the teeth to appear "ghostly" on radiographs. Diastema is a condition in which there is a gap between two teeth caused by the imbalance in the relationship between the jaw and the size of teeth.	Biology and health sciences	Human anatomy	null
68819	https://en.wikipedia.org/wiki/Naturopathy	Naturopathy	Naturopathy, or naturopathic medicine, is a form of alternative medicine. A wide array of practices branded as "natural", "non-invasive", or promoting "self-healing" are employed by its practitioners, who are known as naturopaths. Difficult to generalize, these treatments range from the pseudoscientific and thoroughly discredited, like homeopathy, to the widely accepted, like certain forms of psychotherapy. The ideology and methods of naturopathy are based on vitalism and folk medicine rather than evidence-based medicine, although practitioners may use techniques supported by evidence. The ethics of naturopathy have been called into question by medical professionals and its practice has been characterized as quackery. Naturopathic practitioners commonly encourage alternative treatments that are rejected by conventional medicine, including resistance to surgery or vaccines for some patients. The diagnoses made by naturopaths often have no basis in science and are often not accepted by mainstream medicine. Naturopaths frequently campaign for legal recognition in the United States. Naturopathy is prohibited in three U.S. states (Florida, South Carolina, and Tennessee) and tightly regulated in many others. Some states, however, allow naturopaths to perform minor surgery or even prescribe drugs. While some schools exist for naturopaths, and some jurisdictions allow such practitioners to call themselves doctors, the lack of accreditation, scientific medical training, and quantifiable positive results means they lack the competency of true medical doctors. History The term "naturopathy" originates from "natura" (Latin root for birth) and "pathos" (the Greek root for suffering) to suggest "natural healing". Naturopaths claim the ancient Greek "Father of Medicine", Hippocrates, as the first advocate of naturopathic medicine, before the term existed. Naturopathy has its roots in the 19th-century Natural Cure movement of Europe. In Scotland, Thomas Allinson started advocating his "Hygienic Medicine" in the 1880s, promoting a natural diet and exercise with avoidance of tobacco and overwork. The term naturopathy was coined in 1895 by John Scheel, and purchased by Benedict Lust, whom naturopaths consider to be the "Father of U.S. Naturopathy". Lust had been schooled in hydrotherapy and other natural health practices in Germany by Father Sebastian Kneipp; Kneipp sent Lust to the United States to spread his drugless methods. Lust defined naturopathy as a broad discipline rather than a particular method, and included such techniques as hydrotherapy, herbal medicine, and homeopathy, as well as eliminating overeating, tea, coffee, and alcohol. He described the body in spiritual and vitalistic terms with "absolute reliance upon the cosmic forces of man's nature". According to the Merriam-Webster Dictionary, the first known use of "naturopathy" in print is from 1901. From 1901, Lust founded the American School of Naturopathy in New York. In 1902, the original North American Kneipp Societies were discontinued and renamed "Naturopathic Societies". In September 1919, the Naturopathic Society of America was dissolved and Benedict Lust founded the American Naturopathic Association to supplant it. Naturopaths became licensed under naturopathic or drugless practitioner laws in 25 states in the first three decades of the twentieth century. Naturopathy was adopted by many chiropractors, and several schools offered both Doctor of Naturopathy (ND) and Doctor of Chiropractic (DC) degrees. Estimates of the number of naturopathic schools active in the United States during this period vary from about one to two dozen. After a period of rapid growth, naturopathy went into decline for several decades after the 1930s. In 1910, the Carnegie Foundation for the Advancement of Teaching published the Flexner Report, which criticized many aspects of medical education, especially quality and lack of scientific rigour. The advent of penicillin and other "miracle drugs" and the consequent popularity of modern medicine also contributed to naturopathy's decline. In the 1940s and 1950s, a broadening in scope of practice laws led many chiropractic schools to drop their ND degrees, though many chiropractors continued to practice naturopathy. From 1940 to 1963, the American Medical Association campaigned against heterodox medical systems. By 1958, practice of naturopathy was licensed in only five states. In 1968, the United States Department of Health, Education, and Welfare issued a report on naturopathy concluding that naturopathy was not grounded in medical science and that naturopathic education was inadequate to prepare graduates to make appropriate diagnosis and provide treatment; the report recommends against expanding Medicare coverage to include naturopathic treatments. In 1977 an Australian committee of inquiry reached similar conclusions; it did not recommend licensure for naturopaths. Beginning in the 1970s, there was a revival of interest in the United States and Canada, in conjunction with the "holistic health" movement. , fifteen U.S. states, Puerto Rico, the US Virgin Islands and the District of Columbia licensed naturopathic doctors, and the State of Washington requires insurance companies to offer reimbursement for services provided by naturopathic physicians. On the other hand, some states such as South Carolina and Tennessee prohibit the practice of naturopathy. In the United States, the Indian Health Service began accepting naturopathic doctors in their clinics and practice in 2013, also making loan repayment available to ND's. In 2015, a former naturopathic doctor, Britt Marie Hermes, began writing critically about her experience being trained in and practicing naturopathic medicine. Her blog garnered a large following among skeptics while enraging some proponents of alternative medicine. Practice In 2003, a report was presented by Kimball C. Atwood, an American medical doctor and researcher from Newton, Massachusetts, best known as a critic of naturopathic medicine, stating among other criticisms that "The practice of naturopathy is based on a belief in the body's ability to heal itself through a special vital energy or force guiding bodily processes internally". Diagnosis and treatment concern primarily alternative therapies and "natural" methods that naturopaths claim promote the body's natural ability to heal. Many naturopaths in India now use modern diagnostic techniques in their practice. Naturopaths focus on a holistic approach, avoiding the use of surgery and conventional medicines. Naturopaths aim to prevent illness through stress reduction and changes to diet and lifestyle, often rejecting the methods of evidence-based medicine. A consultation typically begins with a comprehensive patient interview assessing lifestyle, medical history, emotional tone, and physical features, as well as physical examination. Many naturopaths present themselves as primary care providers, and some naturopathic physicians may prescribe drugs, perform minor surgery, and integrate other conventional medical approaches such as diet and lifestyle counselling with their naturopathic practice. Traditional naturopaths deal exclusively with lifestyle changes, not diagnosing or treating disease. Naturopaths do not generally recommend vaccines and antibiotics, based in part on the early views that shaped the profession, and they may provide alternative remedies even in cases where evidence-based medicine has been shown effective. Methods Naturopaths are often opposed to mainstream medicine and take an antivaccinationist stance. The particular modalities used by a naturopath vary with training and scope of practice. These may include herbalism, homeopathy, acupuncture, nature cures, physical medicine, applied kinesiology, colonic enemas, chelation therapy, color therapy, cranial osteopathy, hair analysis, iridology, live blood analysis, ozone therapy, psychotherapy, public health measures and hygiene, reflexology, rolfing, massage therapy, and traditional Chinese medicine. Nature cures include a range of therapies based on exposure to natural elements such as sunshine, fresh air, or heat or cold, as well as nutrition advice such as following a vegetarian and whole food diet, fasting, or abstention from alcohol and sugar. Physical medicine includes naturopathic, osseous, or soft tissue manipulative therapy, sports medicine, exercise, and hydrotherapy. Psychological counseling includes meditation, relaxation, and other methods of stress management. A 2004 survey determined the most commonly prescribed naturopathic therapeutics in Washington state and Connecticut were botanical medicines, vitamins, minerals, homeopathy, and allergy treatments. An examination published in 2011 of naturopathic clinic websites in Alberta and British Columbia found that the most commonly advertised therapies were homeopathy, botanical medicine, nutrition, acupuncture, lifestyle counseling, and detoxification. In 2020, a survey of methods used by naturopaths in fourteen countries reported that 27% of clients received acupuncture, 22% homeopathy, 16% "other energetic medicines", and 13.5% were given hydrotherapy. A mean of 4.0 "treatments" were provided to each customer. One-third (33%) of patients consulted with only the naturopath to manage their primary health concern. Evidence basis Naturopathy as a whole lacks an adequate scientific basis, and it is rejected by the medical community. Although it includes valid lifestyle advice from mainstream medicine (healthy sleep, balanced diet, regular exercise), it typically adds a range of pseudoscientific beliefs. Some methods rely on immaterial "vital energy fields", the existence of which has not been proven, and there is concern that naturopathy as a field tends towards isolation from general scientific discourse. Naturopathy is criticized for its reliance on and its association with unproven, disproven, and other controversial alternative medical treatments, and for its vitalistic underpinnings. Natural substances known as nutraceuticals show little promise in treating diseases, especially cancer, as laboratory experiments have shown limited therapeutic effect on biochemical pathways, while clinical trials demonstrate poor bioavailability. According to the American Cancer Society, "scientific evidence does not support claims that naturopathic medicine can cure cancer or any other disease". According to Britt Hermes, naturopath student programs are problematic because "As a naturopath [student], you are making justifications to make the rules and to fudge the standards of how to interpret research all along the way. Because if you don't, you're not left with anything, basically". In 2015, the Australian Government's Department of Health published the results of a review of alternative therapies that sought to determine if any were suitable for being covered by health insurance; Naturopathy was one of 17 therapies evaluated for which no clear evidence of effectiveness was found. Kimball C. Atwood IV writes, in the journal Medscape General Medicine, In another article, Atwood writes that "Physicians who consider naturopaths to be their colleagues thus find themselves in opposition to one of the fundamental ethical precepts of modern medicine. If naturopaths are not to be judged "nonscientific practitioners", the term has no useful meaning". A former licensed naturopathic doctor, Britt Marie Hermes, states that "any product that is sold by a naturopath almost guarantees that there is no reliable scientific data to support whatever health claims are made, and that while some naturopaths claim to only practice evidence based medicine, "the problem is, all naturopaths in an accredited naturopathic program are required to extensively study homeopathy, herbal medicine, energy healing, chiropractic techniques, water therapy" and other pseudoscientific practices. Hermes further notes that, while some naturopaths claim that their method can be effective treatments for psychological disorders, "no naturopathic treatment has been clinically proven to be safe and effective for bipolar disorder or any other condition." According to Arnold S. Relman, the Textbook of Natural Medicine is inadequate as a teaching tool, as it omits to mention or treat in detail many common ailments, improperly emphasizes treatments "not likely to be effective" over those that are, and promotes unproven herbal remedies at the expense of pharmaceuticals. He concludes that "the risks to many sick patients seeking care from the average naturopathic practitioner would far outweigh any possible benefits". The Massachusetts Medical Society states, "Naturopathic practices are unchanged by research and remain a large assortment of erroneous and potentially dangerous claims mixed with a sprinkling of non-controversial dietary and lifestyle advice." Safety of natural treatments Naturopaths often recommend exposure to naturally occurring substances, such as sunshine, herbs and certain foods, as well as activities they describe as natural, such as exercise, meditation and relaxation. Naturopaths claim that these natural treatments help restore the body's innate ability to heal itself without the adverse effects of conventional medicine. However, "natural" methods and chemicals are not necessarily safer or more effective than "artificial" or "synthetic" ones, and any treatment capable of eliciting an effect may also have deleterious side effects. Certain naturopathic treatments offered by naturopaths, such as homeopathy, rolfing, and iridology, are widely considered pseudoscience or quackery. Stephen Barrett of QuackWatch and the National Council Against Health Fraud has stated that naturopathy is "simplistic and that its practices are riddled with quackery". "Non-scientific health care practitioners, including naturopaths, use unscientific methods and deception on a public who, lacking in-depth health care knowledge, must rely upon the assurance of providers. Quackery not only harms people, it undermines the ability to conduct scientific research and should be opposed by scientists", says William T. Jarvis. In the 2018 Australian case against Marlyin Bodnar, who advised a mother to treat her infant son's eczema with a raw food diet which nearly led to the child's starvation death, Judge Peter Berman said, "Well intentioned but seriously misguided advice is, as the facts of this case demonstrate, capable of causing great harm and even death to vulnerable children." Furthermore, Britt Hermes criticizes the "pervasive culture of patient blaming" among naturopathic practitioners, where "when something doesn't work for the patient and the patient is not experiencing all of the positive effects and zero side-effects that are promised with the therapy, it's never because the therapy doesn't work, it's because the patient didn't do something right." Vaccination Many naturopathy practitioners voice their opposition to vaccination. The reasons for this opposition are based, in part, on the early views which shaped the foundation of this occupation. A naturopathy textbook, co-authored by Joseph Pizzorno, recalls anti-vaccine beliefs associated with the founding of naturopathy in the United States: "a return to nature in regulating the diet, breathing, exercising, bathing and the employment of various forces" in lieu of the smallpox vaccine. In general, evidence about associations between naturopathy and pediatric vaccination is sparse, but "published reports suggest that only a minority of naturopathic physicians actively support full vaccination". In Washington state from 2000 to 2003, children were significantly less likely to receive immunizations if they had seen a naturopath. A survey of naturopathic students published in 2004 found that students at the Canadian College of Naturopathic Medicine became less likely to recommend vaccinations to their patients and became more distrustful of public health and conventional medicine as they advanced in the program. The British Columbia Naturopathic Association lists several major concerns regarding the pediatric vaccine schedule and vaccines in general, and the group's policy is to not advocate for or against vaccines. The Oregon Association of Naturopathic Physicians reports that many naturopaths "customize" the pediatric vaccine schedule. As of April 25, 2022, a British Columbia government report found that 69.2% of naturopaths reported having received at least two COVID vaccines or receiving a medical exemption. This was much lower than all the other regulated medical professions in the report. The number for two professionsdieticians and physicians/surgeonswas 98%. , the American Association of Naturopathic Physicians, which is the largest professional organization for licensed naturopaths in the U.S., is "still discussing its stance on vaccinations". Practitioners Naturopath practitioners can generally be categorized into three groups: 1) those with a government issued license; 2) those who practice outside of an official status ("traditional naturopaths"); 3) those who are primarily another kind of health professional who also practices naturopathy. In Switzerland, these divisions fall between those with a federal diploma, those recognized by health insurances, and those with neither federal diploma nor recognition by health insurances. Naturopaths with federal diploma can be divided into four categories: European traditional medicine, Chinese traditional medicine, ayurvedic medicine and homeopathy. The number of listed naturopaths (including traditional healers) in Switzerland rose from 223 in 1970 to 1835 in 2000. Licensed naturopaths Licensed naturopaths may be referred to as "naturopathic doctors" or "naturopathic physicians" in 26 US states or territories and 5 Canadian provinces. Licensed naturopaths present themselves as primary care providers. Licensed naturopaths do not receive comparable training to medical doctors in terms of the quality of education or quantity of hours. In British Columbia, legislation permits licensed naturopaths to use the title "doctor" or "physician". However, section 102 of the bylaw of the College of Naturopathic Physicians of British Columbia (CNPBC), the terms "naturopathic" or "naturopathic medicine" must be included anytime the term doctor or physician is used by a member of the CNPBC. Education Licensed naturopaths must pass the Naturopathic Physicians Licensing Examinations (NPLEX) administered by the North American Board of Naturopathic Examiners (NABNE) after graduating from a program accredited by the Council on Naturopathic Medical Education (CNME). Training in CNME-accredited programs includes basic medical diagnostics and procedures such as rudimentary physical exams and common blood tests, in addition to pseudoscientific modalities, such as homeopathy, acupuncture, and energy modalities. These accredited programs have been criticized for misrepresenting their medical rigor and teaching subjects that are antithetical to the best understandings of science and medicine. The CNME as an accrediting authority has been characterized as unreliable and suffering from conflicts of interest. The naturopathic licensing exam has been called a mystery by those outside the naturopathic profession and criticized for testing on homeopathic remedies, including for the use to treat pediatric emergencies. Several schools in North America exist for the study of naturopathic medicine, some accredited by the CNME. The CNME and the Association of Accredited Naturopathic Medical Colleges (AANMC) claim entrance requirements and curricula at accredited colleges are often similar or comparable to those required and offered at conventional medical schools. However, the lack of accreditation by the Liaison Committee on Medical Education may indicate insufficiency of scientific medical training and/or quantifiable positive results, and accordingly it remains disputed whether graduates of medical colleges accredited by the CNME have the competency of Medical Doctors and Doctors of Osteopathy. Naturopathic doctors are not eligible for medical residencies, which are available exclusively for medical doctors and doctors of osteopathic medicine. There are limited post-graduate "residency" positions available to naturopathic doctors offered through naturopathic schools and naturopathic clinics approved by the CNME. Most naturopathic doctors do not complete such a residency, and naturopathic doctors are not mandated to complete one for licensure, except in the states of Utah and Connecticut. Continuing education in naturopathic modalities for health care professionals varies greatly. Political activity in the United States Naturopathic practitioners affiliated with the CNME-accredited schools lobby state, provincial, and federal governments for medical licensure and participation in social health programs. The American Association of Naturopathic Physicians represents licensed naturopaths in the United States; the Canadian Association of Naturopathic Doctors represents licensed naturopaths in Canada. Naturopathic lobbying efforts are funded by vitamin and supplement makers and focus on portraying naturopathic education as comparable to medical education received by physicians and on having high professional standards. Medical societies and advocacy groups dispute these claims by citing evidence of licensed naturopathic practitioners using pseudoscientific methods without a sound evidence basis and lacking adequate clinical training to diagnose and treat disease competently according to the standard of care. Jann Bellamy has characterized the process by which naturopathic practitioners and other practitioners of pseudoscience convince lawmakers to provide them with medical licenses as "legislative alchemy". Since 2005, the Massachusetts Medical Society has opposed licensure based on concerns that NDs are not required to participate in residency and concerns that the practices of naturopaths included many "erroneous and potentially dangerous claims". The Massachusetts Special Commission on Complementary and Alternative Medical Practitioners rejected their concerns and recommended licensure. The Massachusetts Medical Society states: In 2015, a former naturopathic doctor, Britt Marie Hermes, who graduated from Bastyr University and practiced as a licensed ND in Washington and Arizona, began advocating against naturopathic medicine. In addition to opposing further licensure, she believes that NDs should not be allowed to use the titles "doctor" or "physician", and be barred from treating children. She states: Traditional naturopaths Traditional naturopaths are represented in the United States by the American Naturopathic Association (ANA), representing about 1,800 practitioners and the American Naturopathic Medical Association (ANMA). The level of naturopathic training varies among traditional naturopaths in the United States. Traditional naturopaths may complete non-degree certificate programs or undergraduate degree programs and generally refer to themselves as naturopathic consultants. These programs often offer online unaccredited degrees, but do not offer comprehensive biomedical education or clinical training. Traditional naturopathic practitioners surveyed in Australia perceive evidence-based medicine to be an ideological assault on their beliefs in vitalistic and holistic principles. They advocate for the integrity of natural medicine practice. Naturopaths graduating from accredited programs argued in 2002 that their training used evidence-based scientific principles unlike traditional naturopathic programs, but this claim remains inaccurate. Regulation Naturopathy is practiced in many countries and is subject to different standards of regulation and levels of acceptance. The scope of practice varies widely between jurisdictions, with some covering naturopathy under medical regulation and allowing practitioners to prescribe drugs and perform minor surgery, while other jurisdictions outlaw naturopathy entirely. Australia In 1977, a Commonwealth Government inquiry reviewed all colleges of naturopathy in Australia and found that despite having syllabuses appearing to cover the basic biomedical sciences, actual lectures had little connection to those syllabuses and no significant practical work was available. In addition, there did not appear to be significant or systematic coverage of techniques favoured by naturopaths, such as homeopathy, Bach's floral remedies, or mineral salts. The position of the Australian Medical Association is that "evidence-based aspects of complementary medicine can be part of patient care by a medical practitioner", but it has concerns that there is "limited efficacy evidence regarding most complementary medicine. Unproven complementary medicines and therapies can pose a risk to patient health either directly through misuse or indirectly if a patient defers seeking medical advice." The AMA's position on regulation is that "there should be appropriate regulation of complementary medicine practitioners and their activities". In 2015, the Australian government found no clear evidence of effectiveness for naturopathy. Accordingly, In 2017 the Australian government named naturopathy as a practice that would not qualify for insurance subsidies, saying this step would "ensure taxpayer funds are expended appropriately and not directed to therapies lacking evidence". India In India, naturopathy is overseen by the Department of Ayurveda, Yoga and Naturopathy, Unani, Siddha and Homoeopathy (AYUSH); there is a 5½-year degree in "Bachelor of Naturopathy and Yogic Sciences" (BNYS) degree that was offered by twelve colleges in India . The National Institute of Naturopathy in Pune that operates under AYUSH, which was established on December 22, 1986, and encourages facilities for standardization and propagation of the existing knowledge and its application through research in naturopathy throughout India. North America In five Canadian provinces, seventeen U.S. states, and the District of Columbia, naturopathic doctors who are trained at an accredited school of naturopathic medicine in North America are entitled to use the designation ND or NMD. Elsewhere, the designations "naturopath", "naturopathic doctor", and "doctor of natural medicine" are generally unprotected or prohibited. In North America, each jurisdiction that regulates naturopathy defines a local scope of practice for naturopathic doctors that can vary considerably. Some regions permit minor surgery, access to prescription drugs, spinal manipulations, midwifery (natural childbirth), and gynecology; other regions exclude these from the naturopathic scope of practice or prohibit the practice of naturopathy entirely. Canada Five Canadian provinces license naturopathic doctors: Ontario, British Columbia, Manitoba, Saskatchewan, and Alberta. British Columbia has the largest scope of practice in Canada, allowing certified NDs to prescribe pharmaceuticals and perform minor surgeries. Ontario also permits prescription from a modified formulary list, following separate examination. United States U.S. jurisdictions that currently regulate or license naturopathy include Alaska, Arizona, California, Connecticut, Colorado, Delaware, District of Columbia, Hawaii, Idaho, Kansas, Maine, Maryland, Massachusetts, Minnesota, Montana, New Hampshire, New Mexico, North Dakota, Oregon, Pennsylvania, Rhode Island, Wisconsin, Puerto Rico, US Virgin Islands, Utah, Vermont, and Washington. Additionally, Virginia licenses the practice of naturopathy under a grandfather clause. (This was previously also the case in Florida, though currently no practitioners remain active under the grandfather provisions). U.S. jurisdictions that permit access to prescription drugs: Arizona, California, District of Columbia, Hawaii, Idaho, Kansas, Maine, Montana, New Hampshire, New Mexico, Oregon, Utah, Vermont, and Washington. U.S. jurisdictions that permit minor surgery: Arizona, District of Columbia, Kansas, Maine, Montana, Oregon, Utah, Vermont, and Washington. Three U.S. states specifically prohibit the practice of naturopathy: Florida, South Carolina and Tennessee. Switzerland The Swiss Federal Constitution defines the Swiss Confederation and the Cantons of Switzerland within the scope of their powers to oversee complementary medicine. In particular, the Federal authorities must set up diplomas for the practice of non-scientific medicine. The first of such diplomas has been validated in April 2015 for the practice of naturopathy. There is a long tradition of naturopathy and traditional medicine in Switzerland. The Cantons of Switzerland make their own public health regulations. Although the law in certain cantons is typically monopolistic, the authorities are relatively tolerant with regard to alternative practitioners. United Kingdom Naturopathy is not regulated in the United Kingdom. In 2012, publicly funded universities in the United Kingdom dropped their alternative medicine programs, including naturopathy.	Biology and health sciences	Alternative and traditional medicine	Health
68833	https://en.wikipedia.org/wiki/Iteration	Iteration	Iteration is the repetition of a process in order to generate a (possibly unbounded) sequence of outcomes. Each repetition of the process is a single iteration, and the outcome of each iteration is then the starting point of the next iteration. In mathematics and computer science, iteration (along with the related technique of recursion) is a standard element of algorithms. Mathematics In mathematics, iteration may refer to the process of iterating a function, i.e. applying a function repeatedly, using the output from one iteration as the input to the next. Iteration of apparently simple functions can produce complex behaviors and difficult problems – for examples, see the Collatz conjecture and juggler sequences. Another use of iteration in mathematics is in iterative methods which are used to produce approximate numerical solutions to certain mathematical problems. Newton's method is an example of an iterative method. Manual calculation of a number's square root is a common use and a well-known example. Computing In computing, iteration is the technique marking out of a block of statements within a computer program for a defined number of repetitions. That block of statements is said to be iterated; a computer scientist might also refer to that block of statements as an "iteration". Implementations Loops constitute the most common language constructs for performing iterations. The following pseudocode "iterates" three times the line of code between begin & end through a for loop, and uses the values of i as increments. a := 0 for i := 1 to 3 do { loop three times } begin a := a + i; { add the current value of i to a } end; print(a); { the number 6 is printed (0 + 1; 1 + 2; 3 + 3) } It is permissible, and often necessary, to use values from other parts of the program outside the bracketed block of statements, to perform the desired function. Iterators constitute alternative language constructs to loops, which ensure consistent iterations over specific data structures. They can eventually save time and effort in later coding attempts. In particular, an iterator allows one to repeat the same kind of operation at each node of such a data structure, often in some pre-defined order. Iteratees are purely functional language constructs, which accept or reject data during the iterations. Relation with recursion Recursions and iterations have different algorithmic definitions, even though they can generate identical effects/results. The primary difference is that recursion can be employed as a solution without prior knowledge as to how many times the action will have to repeat, while a successful iteration requires that foreknowledge. Some types of programming languages, known as functional programming languages, are designed such that they do not set up a block of statements for explicit repetition, as with the for loop. Instead, those programming languages exclusively use recursion. Rather than call out a block of code to be repeated a pre-defined number of times, the executing code block instead "divides" the work to be done into a number of separate pieces, after which the code block executes itself on each individual piece. Each piece of work will be divided repeatedly until the "amount" of work is as small as it can possibly be, at which point the algorithm will do that work very quickly. The algorithm then "reverses" and reassembles the pieces into a complete whole. The classic example of recursion is in list-sorting algorithms, such as merge sort. The merge sort recursive algorithm will first repeatedly divide the list into consecutive pairs; each pair is then ordered, then each consecutive pair of pairs, and so forth until the elements of the list are in the desired order. The code below is an example of a recursive algorithm in the Scheme programming language that will output the same result as the pseudocode under the previous heading. (let iterate ((i 1) (a 0)) (if (<= i 3) (iterate (+ i 1) (+ a i)) (display a))) Education In some schools of pedagogy, iterations are used to describe the process of teaching or guiding students to repeat experiments, assessments, or projects, until more accurate results are found, or the student has mastered the technical skill. This idea is found in the old adage, "Practice makes perfect." In particular, "iterative" is defined as the "process of learning and development that involves cyclical inquiry, enabling multiple opportunities for people to revisit ideas and critically reflect on their implication." Unlike computing and math, educational iterations are not predetermined; instead, the task is repeated until success according to some external criteria (often a test) is achieved.	Technology	Software development: General	null
68853	https://en.wikipedia.org/wiki/Dahlia	Dahlia	Dahlia ( , ) is a genus of bushy, tuberous, herbaceous perennial plants native to Mexico and Central America. Dahlias are members of the Asteraceae (synonym name: Compositae) family of dicotyledonous plants, its relatives include the sunflower, daisy, chrysanthemum, and zinnia. There are 49 species of dahlia, with flowers in almost every hue (except blue), with hybrids commonly grown as garden plants. Dahlias were known only to the Aztecs and other southern North American peoples until the Spanish conquest, after which the plants were brought to Europe. The tubers of some varieties are of medicinal and dietary value to humans because they contain inulin, a polymer of the fruit sugar, fructose. Description Dahlias are perennial plants with tuberous roots. They are not frost hardy, and require protection from frost if grown in regions with cold winters. While some have herbaceous stems, others have stems which lignify in the absence of secondary tissue and resprout following winter dormancy, allowing further seasons of growth. As members of the family Asteraceae, dahlias have composite flower heads called capitula that are composed of multiple florets arranged in a central disc with surrounding petal-like rays. Each floret is a flower in its own right. The modern name Asteraceae comes from the type genus Aster and the Ancient Greek word for "star", referring to the appearance of a star with surrounding rays. The stems are leafy, ranging in height from as low as to more than . Flower forms are variable, with one head per stem; these can be as small as in diameter or up to ("dinner plate"). The majority of species do not produce scented flowers. Like most plants that do not attract pollinating insects through scent, they are brightly colored, displaying most hues, with the exception of blue. The great variety in species results from garden dahlias being octoploids, having eight sets of homologous chromosomes. In addition, dahlias also contain many transposons—genetic pieces that move from place to place upon an allele—which contributes to their manifesting such great diversity. Taxonomy Taxonomic history Early history Spaniards reported finding the plants growing in Mexico in 1525, but the earliest known description is by Francisco Hernández, physician to Philip II, who was ordered to visit Mexico in 1570 to study the "natural products of that country". They were used as a source of food by the indigenous peoples, who both gathered wild specimens and cultivated crops. The indigenous peoples variously identified the plants as "Chichipatl" (Toltecs) and "Acocotle" or "Cocoxochitl" (Aztecs). From Hernandez's perception of Nahuatl to Spanish (through various other translations) the word is "water cane", "water pipe", "water pipe flower", "hollow stem flower", or "cane flower", all referring to the hollow plant stems. Hernandez described two varieties of dahlias (the pinwheel-like Dahlia pinnata and the huge Dahlia imperialis) as well as other medicinal plants of New Spain. Francisco Dominguez, an Hidalgo gentleman who accompanied Hernandez on part of his seven-year study, made a series of drawings to supplement the four volume report. Three of his drawings showed plants with flowers: two resembled the modern bedding dahlia, and one resembled the species Dahlia merckii; all displayed a high degree of doubleness. In 1578, a manuscript titled Nova Plantarum, Animalium et Mineralium Mexicanorum Historia, was sent back to the Escorial in Madrid. It was translated into Latin by Francisco Ximenes in 1615. In 1640, Francisco Cesi, President of the Academia dei Lincei of Rome, bought the Ximenes translation and, after annotating it, published it in 1649–1651 as two volumes, Rerum Medicarum Novae Hispaniae Thesaurus Seu Nova Plantarium, Animalium et Mineralium Mexicanorum Historia. The original manuscripts were destroyed in a fire in the mid-1600s. European introduction In 1787, the French botanist Nicolas-Joseph Thiéry de Menonville, sent to Mexico to steal the cochineal insect valued for its scarlet dye, reported the strangely beautiful flowers he had seen growing in a garden in Oaxaca. In 1789, Vicente Cervantes, director of the Botanical Garden at Mexico City, sent "plant parts" to Abbe Antonio José Cavanilles, director of the Royal Gardens of Madrid. Cavanilles flowered one plant that same year, then the second one a year later. In 1791 he called the new growths "Dahlia" for Anders Dahl. The first plant was called Dahlia pinnata after its pinnate foliage; the second, Dahlia rosea for its rose-purple color. In 1796, from the parts sent by Cervantes, Cavanilles flowered a third plant, which he named Dahlia coccinea for its scarlet color. In 1798, Cavanilles sent D. pinnata seeds to Parma, Italy. That year, the Marchioness of Bute, wife of the Earl of Bute, the English Ambassador to Spain, obtained a few seeds from Cavanilles and sent them to Kew Gardens, where they flowered but were lost after two to three years. In the following years Madrid sent seeds to Berlin and Dresden in Germany, and to Turin and Thiene in Italy. In 1802, Cavanilles sent tubers of "these three" (D. pinnata, D. rosea, D. coccinea) to Swiss botanist Augustin Pyramus de Candolle at University of Montpelier in France, Andre Thouin at the Jardin des Plantes in Paris and Scottish botanist William Aiton at Kew Gardens. That same year, John Fraser, English nurseryman and later botanical collector to the Czar of Russia, brought D. coccinea seeds from Paris to the Apothecaries Gardens in England, where they flowered in his greenhouse a year later, providing Botanical Magazine with an illustration. In 1804, a new species, Dahlia sambucifolia, was successfully grown at Holland House, Kensington. Whilst in Madrid in 1804, Lady Holland was given either dahlia seeds or tubers by Cavanilles. She sent them back to England, to Lord Holland's librarian at Holland House, who successfully raised the plants and produced two double flowers a year later. The plants raised in 1804 did not survive; new stock was brought from France in 1815. In 1824, Lord Holland sent his wife a note containing the following verse:The dahlia you brought to our isle Your praises for ever shall speak; Mid gardens as sweet as your smile, And in colour as bright as your cheek. In 1805, German naturalist Alexander von Humboldt sent more seeds from Mexico to Aiton in England, Thouin in Paris, and Christoph Friedrich Otto, director of the Berlin Botanical Garden. More significantly, he sent seeds to botanist Carl Ludwig Willdenow in Germany. Willdenow now reclassified the rapidly growing number of species, changing the genus from Dahlia to Georgina; after naturalist Johann Gottlieb Georgi. He combined the Cavanilles species D. pinnata and D. rosea under the name of Georgina variabilis; D. coccinea was still held to be a separate species, which he renamed Georgina coccinea. Classification Since 1789 when Cavanilles first flowered the dahlia in Europe, there has been an ongoing effort by many growers, botanists and taxonomists, to determine the development of the dahlia to modern times. At least 85 species have been reported: approximately 25 of these were first reported from the wild; the remainder appeared in gardens in Europe. They were considered hybrids, the results of crossing between previously reported species, or developed from the seeds sent by Humboldt from Mexico in 1805, or perhaps from some other undocumented seeds that had found their way to Europe. Several of these were soon discovered to be identical with earlier reported species, but the greatest number are new varieties. Morphological variation is highly pronounced in the dahlia. William John Cooper Lawrence, who hybridized hundreds of families of dahlias in the 1920s, stated: "I have not yet seen any two plants in the families I have raised which were not to be distinguished one from the other. Constant reclassification of the 85 reported species has resulted in a considerably smaller number of distinct species, as there is a great deal of disagreement today between systematists over classification. In 1829, all species growing in Europe were reclassified under an all-encompassing name of D. variabilis, Desf., though this is not an accepted name. Through the interspecies cross of the Humboldt seeds and the Cavanilles species, 22 new species were reported by that year, all of which had been classified in different ways by several different taxonomists, creating considerable confusion as to which species was which. As of now Dahlias are classified into 15 different species by botanist Liberty Hyde Bailey. In 1830 William Smith suggested that all dahlia species could be divided into two groups for color, red-tinged and purple-tinged. In investigating this idea Lawrence determined that with the exception of D. variabilis, all dahlia species may be assigned to one of two groups for flower-colour: Group I (ivory-magenta) or Group II (yellow-orange-scarlet). Modern classification The genus Dahlia is situated in the Asteroideae subfamily of the Asteraceae, in the Coreopsideae tribe. Within that tribe it is the second largest genus, after Coreopsis, and appears as a well defined clade within the Coreopsideae. Subdivision Infrageneric subdivision Sherff (1955), in the first modern taxonomy described three sections for the 18 species he recognised, Pseudodendron, Epiphytum and Dahlia. By 1969 Sørensen recognised 29 species and four sections by splitting off Entemophyllon from section Dahlia. By contrast Giannasi (1975) using a phytochemical analysis based on flavonoids, reduced the genus to just two sections, Entemophyllon and Dahlia, the latter having three subsections, Pseudodendron, Dahlia, and Merckii. Sørensen then issued a further revision in 1980, incorporating subsection Merckii in his original section Dahlia. When he described two new species in the 1980s (Dahlia tubulata and D. congestifolia), he placed them within his existing sections. A further species, Dahlia sorensenii was added by Hansen and Hjerting in (1996). At the same time they demonstrated that Dahlia pinnata should more properly be designated D. x pinnata. D. x pinnata was shown to actually be a variant of D. sorensenii that had acquired hybrid qualities before it was introduced to Europe in the sixteenth century and formally named by Cavanilles. The original wild D. pinnata is presumed extinct. Further species continue to be described, Saar (2003) describing 35 species. However separation of the sections on morphological, cytologal and biocemical criteria has not been entirely satisfactory. To date these sectional divisions have not been fully supported phylogenetically, which demonstrate only section Entemophyllon as a distinct sectional clade. The other major grouping is the core Dahlia clade (CDC), which includes most of the section Dahlia. The remainder of the species occupy what has been described as the variable root clade (VRC) which includes the small section Pseudodendron but also the monotypic section Epiphytum and a number of species from within section Dahlia. Outside of these three clades lie D. tubulata and D. merckii as a polytomy. Horticulturally the sections retain some usage, section Pseudodendron being referred to as 'Tree Dahlias', Epiphytum as the 'Vine Dahlia'. The remaining two herbaceous sections being distinguished by their pinnules, opposing (Dahlia) or alternating (Entemophyllon). Sections Sections (including chromosome numbers), with geographical distribution; Epiphytum Sherff (2n = 32) 10 m tall climber with aerial roots 5 cm thick and up to more than 20 m long; pinnules opposite 1 species, D. macdougallii Sherff Mexico: Oaxaca Entemophyllon P. D. Sorensen (2n = 34) 6 species Mexico: Hidalgo, Nuevo León, Tamaulipas, Querétaro, Durango, San Luis Potosí Pseudodendron P. D. Sorensen (2n = 32) 3 species + D. excelsa of uncertain identity Mexico: Chiapas, Guerrero, Jalisco, Michoacan, Oaxaca, and Costa Rica, El Salvador, Guatemala & Colombia Dahlia (2n = 32, 36 or 64) 24 species Mexico: Distrito Federal, Guerrero, Hidalgo, Morelos, Nuevo León, Puebla, San Luis Potosí, Tamaulipas, Veracruz, Oaxaca, Puebla, Chiapas, México, Huehuetenango, Chihuahua, Durango, Michoacan & Guatemala Only Pseudodendron (D. imperialis) and Dahlia (D. australis, D. coccinea) occur outside Mexico. Species There are currently 42 accepted species in the genus Dahlia but new species continue to be described. Etymology The naming of the plant itself has long been a subject of some confusion. Many sources state that the name "Dahlia" was bestowed by the pioneering Swedish botanist and taxonomist Carl Linnaeus to honor his late student, Anders Dahl, author of Observationes Botanicae. However, Linnaeus died in 1778, more than eleven years before the plant was introduced into Europe in 1789, so while it is generally agreed that the plant was named in 1791 in honor of Dahl, who had died two years before, Linnaeus could not have been the one who did so. It was probably Abbe Antonio Jose Cavanilles, Director of the Royal Gardens of Madrid, who should be credited with the attempt to scientifically define the genus, since he not only received the first specimens from Mexico in 1789, but named the first three species that flowered from the cuttings. Regardless of who bestowed it, the name was not so easily established. In 1805, German botanist Carl Ludwig Willdenow, asserting that the genus Dahlia Thunb. (published a year after Cavanilles's genus and now considered a synonym of Trichocladus) was more widely accepted, changed the plants' genus from Dahlia to Georgina (after the German-born naturalist Johann Gottlieb Georgi, a professor at the Imperial Academy of Sciences of St. Petersburg, Russia). He also reclassified and renamed the first three species grown, and identified, by Cavanilles. It was not until 1810, in a published article, that he officially adopted the Cavanilles's original designation of Dahlia. However, the name Georgina still persisted in Germany for the next few decades. In Russian, it is still named Georgina (). Distribution and habitat Dahlia is found predominantly in Mexico, but some species are found ranging as far south as northern South America. D. australis occurs at least as far south as southwestern Guatemala, while D. coccinea and D. imperialis also occur in parts of Central America and northern South America. Dahlia is a genus of the uplands and mountains, being found at elevations between , in what has been described as a "pine-oak woodland" vegetative zone. Most species have limited ranges scattered throughout many mountain ranges in Mexico Ecology The most common pollinators are bees and small beetles. Pests and diseases Slugs and snails are serious pests in some parts of the world, particularly in spring when new growth is emerging through the soil. Earwigs can also disfigure the blooms and foliage. The other main pests likely to be encountered are aphids (usually on young stems and immature flower buds), red spider mite (causing foliage mottling and discolouration, worse in hot and dry conditions) and capsid bugs (resulting in contortion and holes at growing tips). Diseases affecting dahlias include powdery mildew, grey mould (Botrytis cinerea), verticillium wilt, dahlia smut (Entyloma calendulae f. dahliae), phytophthora and some plant viruses. Dahlias are a source of food for the larvae of some Lepidoptera species including angle shades, common swift, ghost moth and large yellow underwing. Cultivation Dahlias grow naturally in climates that do not experience frost (the tubers are hardy to USDA Zone 8), consequently they are not adapted to withstand sub-zero temperatures. However, their tuberous nature enables them to survive periods of dormancy, and this characteristic means that gardeners in temperate climates with frosts can grow dahlias successfully, provided the tubers are lifted from the ground and stored in cool yet frost-free conditions during the winter. Planting the tubers quite deep () also provides some protection. When in active growth, modern dahlia hybrids perform most successfully in well-watered yet free-draining soils, in situations receiving plenty of sunlight. Taller cultivars usually require some form of staking as they grow, and all garden dahlias need deadheading regularly, once flowering commences. Horticultural classification Horticulturally the garden dahlia is usually treated as the cultigen D. variabilis hort., which while being responsible for thousands of cultivars has an obscure taxonomic status. History The inappropriate term D. variabilis is often used to describe the cultivars of Dahlia since the correct parentage remains obscure, but probably involves Dahlia coccinea. In 1846 the Caledonia Horticultural Society of Edinburgh offered a prize of 2,000 pounds to the first person succeeding in producing a blue dahlia. This has to date not been accomplished. While dahlias produce anthocyanin, an element necessary for the production of the blue, to achieve a true blue color in a plant, the anthocyanin delphinidin needs six hydroxyl groups. To date, dahlias have only developed five, so the closest that breeders have come to achieving a "blue" specimen are variations of mauve, purples and lilac hues. By the beginning of the twentieth century, a number of different types were recognised. These terms were based on shape or colour, and the National Dahlia Society included cactus, pompon, single, show and fancy in its 1904 guide. Many national societies developed their own classification systems until 1962 when the International Horticultural Congress agreed to develop an internationally recognised system at its Brussels meeting that year, and subsequently in Maryland in 1966. This culminated in the 1969 publication of The International Register of Dahlia Names by the Royal Horticultural Society which became the central registering authority. This system depended primarily on the visibility of the central disc, whether it was open-centred or whether only ray florets were apparent centrally (double bloom). The double-bloom cultivars were then subdivided according to the way in which they were folded along their longitudinal axis: flat, involute (curled inwards) or revolute (curling backwards). If the end of the ray floret was split, they were considered fimbriated. Based on these characteristics, nine groups were defined plus a tenth miscellaneous group for any cultivars not fitting the above characteristics. Fimbriated dahlias were added in 2004, and two further groups (Single and Double orchid) in 2007. The last group to be added, Peony, first appeared in 2012. In many cases the bloom diameter was then used to further label certain groups from miniature to giant. This practice was abandoned in 2012. Modern system (RHS) There are now more than 57,000 registered cultivars, which are officially registered through the Royal Horticultural Society (RHS). The official register is The International Register of Dahlia Names 1969 (1995 reprint) which is updated by annual supplements. The original 1969 registry published about 14,000 cultivars adding a further 1700 by 1986 and in 2003 there were 18,000. Since then about a hundred new cultivars are added annually. Flower type The official RHS classification lists fourteen groups, grouped by flower type, together with the abbreviations used by the RHS; Flower size Earlier versions of the registry subdivided some groups by flower size. Groups 4, 5, 8 and 9 were divided into five subgroups (A to E) from Giant to Miniature, and Group 6 into two subgroups, Small and Miniature. Dahlias were then described by Group and Subgroup, e.g. 5(d) ‘Ace Summer Sunset’. Some Dahlia Societies have continued this practice, but this is neither official nor standardised. As of 2013 The RHS uses two size descriptors Dwarf Bedder (Dw.B.) – not usually exceeding in height, e.g. 'Preston Park' (Sin/DwB) Lilliput dahlias (Lil) – not usually exceeding in height, with single, semi-double or double florets up to in diameter. ("baby" or "top-mix" dahlias), e.g. 'Harvest Tiny Tot' (Misc/Lil) Sizes can range from tiny micro dahlias with flowers less than 50 mm to giants that are over 250 mm in diameter. The groupings listed here are from the New Zealand Society: Giant-flowered cultivars have blooms with a diameter over 250 mm. Large-flowered cultivars have blooms with a diameter of 200–250 mm. Medium-flowered cultivars have blooms with a diameter of 155–350 mm. Small-flowered cultivars have blooms with a diameter of 115–155 mm. Miniature-flowered cultivars have blooms with a diameter of 50–115 mm. Pompom-flowered cultivars have blooms with a diameter less than 50 mm. In addition to the official classification and the terminology used by various dahlia societies, individual horticulturalists use a wide range of other descriptions, such as 'Incurved' and abbreviations in their catalogues, such as CO for Collarette. Branding Some plant growers include their brand name in the cultivar name. Thus Fides (part of the Dümmen Orange Group) in the Netherlands developed a series of cultivars which they named the Dahlinova series, for example Dahlinova 'Carolina Burgundy'. These are Group 10 Miscellaneous in the RHS classification scheme. Double dahlias In 1805, several new species were reported with red, purple, lilac, and pale yellow coloring, and the first true double flower was produced in Belgium. One of the more popular concepts of dahlia history, and the basis for many different interpretations and confusion, is that all the original discoveries were single-flowered types, which, through hybridization and selective breeding, produced double forms. Many of the species of dahlias then, and now, have single-flowered blooms. D. coccinea, the third dahlia to bloom in Europe, was a single. But two of the three drawings of dahlias by Dominguez, made in Mexico between 1570 and 1577, showed definite characteristics of doubling. In the early days of the dahlia in Europe, the word "double" simply designated flowers with more than one row of petals. The greatest effort was now directed to developing improved types of double dahlias. During the years 1805 to 1810 several people claimed to have produced a double dahlia. In 1805 Henry C. Andrews made a drawing of such a plant in the collection of Lady Holland, grown from seedlings sent that year from Madrid. Like other doubles of the time it did not resemble the doubles of today. The first modern double, or full double, appeared in Belgium; M. Donckelaar, Director of the Botanic Garden at Louvain, selected plants for that characteristic, and within a few years secured three fully double forms. By 1826 double varieties were being grown almost exclusively, and there was very little interest in the single forms. Up to this time all the so-called double dahlias had been purple, or tinged with purple, and it was doubted if a variety untinged with that color was obtainable. In 1843, scented single forms of dahlias were first reported in Neu Verbass, Austria. D. crocea, a fragrant variety grown from one of the Humboldt seeds, was probably interbred with the single D. coccinea. A new scented species would not be introduced until the next century when the D. coronata was brought from Mexico to Germany in 1907. The exact date the dahlia was introduced in the United States is uncertain. One of the first dahlias in the USA may have been the D. coccinea speciosissima grown by William Leathe, of Cambridgeport, near Boston, around 1929. According to Edward Sayers, "it attracted much admiration, and at that time was considered a very elegant flower, it was however soon eclipsed by that splendid scarlet, the Countess of Liverpool". However, nine cultivars were already listed in the catalog from Thornburn, 1825. And even earlier reference can be found in a catalogue from the Linnaean Botanical Garden, New York, 1820, that includes one scarlet, one purple, and two double orange Dahlias for sale. Sayers stated that "No person has done more for the introduction and advancement of the culture of the Dahlia than George C. Thorburn, of New York, who yearly flowers many thousand plants at his place at Hallet's Cove, near Harlaem. The show there in the flowering season is a rich treat for the lovers of floriculture : for almost every variety can be seen growing in two large blocks or masses which lead from the road to the dwelling-house, and form a complete field of the Dahlia as a foreground to the house. Mr. T. Hogg, William Read, and many other well-known florists have also contributed much in the vicinity of New York, to the introduction of the Dahlia. Indeed so general has become the taste that almost every garden has its show of the Dahlia in the season." In Boston too there were many collections, a collection from the Messrs Hovey of Cambridgeport was also mentioned. In 1835 Thomas Bridgeman, published a list of 160 double dahlias in his "Florist's Guide". 60 of the choicest were supplied by Mr. G. C. Thornburn of Astoria, New York, who got most of them from contacts in the UK. Not a few of them had taken prizes "at the English and American exhibitions". "Stars of the Devil" In 1872 J. T. van der Berg of Utrecht in the Netherlands received a shipment of seeds and plants from a friend in Mexico. The entire shipment was badly rotted and appeared to be ruined, but van der Berg examined it carefully and found a small piece of root that seemed alive. He planted and carefully tended it; it grew into a plant that he identified as a dahlia. He made cuttings from the plant during the winter of 1872–1873. This was an entirely different type of flower, with rich, red color and a high degree of doubling. In 1874 van der Berg catalogued it for sale, calling it Dahlia juarezii to honor Mexican President Benito Pablo Juarez, who had died the year before, and described it as "...equal to the beautiful color of the red poppy. Its form is very outstanding and different in every respect of all known dahlia flowers." This plant has perhaps had a greater influence on the popularity of the modern dahlia than any other. Called "Les Etoiles du Diable" (Stars of the Devil) in France and "Cactus dahlia" elsewhere, the edges of its petals rolled backwards, rather than forward, and this new form revolutionized the dahlia world. It was thought to be a distinct mutation since no other plant that resembled it could be found in the wild. Today it is assumed that D. juarezii had, at one time, existed in Mexico and subsequently disappeared. Nurserymen in Europe crossbred this plant with dahlias discovered earlier; the results became the progenitors of all modern dahlia hybrids today. Award of Garden Merit (RHS) As of 2015, 124 dahlia cultivars have gained the Royal Horticultural Society's Award of Garden Merit, including: "Bednall beauty" "Bishop of Llandaff" "Clair de lune" "David Howard" "Ellen Huston" "Fascination" "Gallery Art Deco" "Gallery Art Nouveau" "Glorie van Heemstede" "Honka" "Moonfire" "Twyning's After Eight" Uses The Aztecs used dahlias to treat epilepsy, and employed the long hollow stem of the D. imperialis for water pipes. Europeans attempted to introduce the tubers as a food crop, but this was unpopular. The dahlia is considered one of the native ingredients in Oaxacan cuisine; several cultivars are still grown especially for their large, sweet potato-like tubers. Dacopa, an intense mocha-tasting extract from the roasted tubers, is used to flavor beverages throughout Central America. Medicine In Europe and America, prior to the discovery of insulin in 1923, diabetics—as well as consumptives—were often given a substance called Atlantic starch or diabetic sugar, derived from inulin, a naturally occurring form of fruit sugar, extracted from dahlia tubers. Inulin is still used in clinical tests for kidney functionality. In culture Founded in 1936, the Bloemencorso Zundert is the largest flower parade in the world entirely made by volunteers using the dahlia. The parade takes place on the first Sunday of September in Zundert, Netherlands. The floats are large artworks made of steel wire, cardboard, papier-mâché and flowers. In the Bloemencorso Zundert, mostly dahlias are used to decorate the objects and it takes thousands of them just to cover one float. Around 8 million dahlias are needed for the entire corso. Of these, around 6 million are cultivated in Zundert. The dahlia was declared the national flower of Mexico in 1963.	Biology and health sciences	Asterales	null
68889	https://en.wikipedia.org/wiki/Biblical%20and%20Talmudic%20units%20of%20measurement	Biblical and Talmudic units of measurement	Biblical and Talmudic units of measurement were used primarily by ancient Israelites and appear frequently within the Hebrew Bible as well as in later rabbinic writings, such as the Mishnah and Talmud. These units of measurement continue to be used in functions regulating Orthodox Jewish contemporary life, based on halacha. The specificity of some of the units used and which are encompassed under these systems of measurement (whether in linear distance, weight or volume of capacity) have given rise, in some instances, to disputes, owing to the discontinuation of their Hebrew names and their replacement by other names in modern usage. Note: The listed measurements of this system range from the lowest to highest acceptable halakhic value, in terms of conversion to and from contemporary systems of measurement. Contemporary unit conversion While documentation on each unit's relation to another's is plentiful, there is much debate, both within Judaism and in academia, about the exact relationship between measurements in the system and those in other measurement systems. Classical definitions, such as that an etzba was seven barleycorns laid side by side, or that a log was equal to six medium-sized eggs, are also open to debate. Nevertheless, the entire system of measurement bears profound resemblance to the Babylonian and the ancient Egyptian systems, and is currently understood to have likely been derived from some combination of the two. Scholars commonly infer the absolute sizes based on the better-known Babylonian units' relations to their contemporary counterparts. Length and distance The original measures of length were clearly derived from the human body—the finger, hand, arm, span, foot, and pace—but since these measures differ between individuals, they are reduced to a certain standard for general use. The Israelite system thus used divisions of the digit or fingerbreadth (Hebrew: אצבע, etzba; plural etzba'ot), the palm or handbreadth (Hebrew: טפח, ; plural /), which is equal to four fingerbreadths, the span (Hebrew: זרת, ), the ell or cubit (Hebrew: אמה, Amah, plural Amot), the mile (Hebrew: מיל, mil; plural milim), and the parsa (Hebrew: פרסה, parasa). The latter two are loan words into the Hebrew language, and borrowed measurements - the Latin mille, and Iranian parasang, respectively; both were units of itinerant distance, and thus varied according to terrain and stride length, and, in the case of the parasang, also on the speed of travel. The Israelite measurements were related as follows: 1 palm [handbreadth] () = 4 digit (etzba'ot) 1 span () = 3 palms () 1 ell [cubit] () = 2 spans (), or 6 palms [handbreadths] 1 mil () = 2000 ells [cubits] () 1 parasang () = 4 mils () Discrepancies of ell The biblical ell is closely related to the cubit, but two different factors are given in the Bible; Ezekiel's measurements imply that the ell was equal to 1 cubit plus 1 palm (Tefah), while elsewhere in the Bible, the ell is equated with 1 cubit exactly. Ezekiel's ell, by which he gave measurements in his guided vision through a future Jerusalem Temple, is thus one sixth larger than the standard ell, for which an explanation seems to be suggested by the Book of Chronicles; the Chronicler writes that Solomon's Temple was built according to "cubits following the first measure", suggesting that over the course of time the original ell was supplanted by a smaller one. The Egyptians also used two different ells, one of which—the royal ell—was a sixth larger than the common ell; this royal measurement was the earlier of the two in Egyptian use, and the one which the Pyramids of the 3rd and 4th Dynasties seem to be measured in integer multiples of. The smaller of the Egyptian ells measured , but the standard Babylonian ell, cast in stone on one of the statues of King Gudea, was 49.5 cm (19.49 in), and the larger Egyptian ell was between 52.5 and 52.8 cm (20.67 and 20.79 in). The Books of Samuel portray the Temple as having a Phoenician architect, and in Phoenicia it was the Babylonian ell which was used to measure the size of parts of ships. Thus scholars are uncertain whether the standard Biblical ell would have been 49.5 or 52.5 cm (19.49 or 20.67 in), but are fairly certain that it was one of these two figures. From these figures for the size of a Biblical ell, that of the basic unit—the finger-breadth (Etzba)—can be calculated to be either 2.1 or 2.2 cm (0.83 or 0.87 in); Rav Avraham Chaim Naeh approximates at 2 cm (0.79 in); Talmudic scholar Chazon Ish at 2.38 cm (0.94 in). The mile (Mil) is thus about 963 or 1146 meters (3160 or 3760 ft)—approximately six or seven tenths of a mile, and significantly shorter than the modern statute or land mile of 5280 ft or 1760 yd (approximately 1.6 km). The precise width of the etzba (finger) has been a subject of controversy among halakhic authorities. The best known are those of the Rav Chayim No'eh and Chazon Ish.	Physical sciences	Measurement systems	Basics and measurement
68954	https://en.wikipedia.org/wiki/Nostril	Nostril	A nostril (or naris , : nares ) is either of the two orifices of the nose. They enable the entry and exit of air and other gasses through the nasal cavities. In birds and mammals, they contain branched bones or cartilages called turbinates, whose function is to warm air on inhalation and remove moisture on exhalation. Fish do not breathe through noses, but they do have two small holes used for smelling, which can also be referred to as nostrils (with the exception of Cyclostomi, which have just one nostril). In humans, the nasal cycle is the normal ultradian cycle of each nostril's blood vessels becoming engorged in swelling, then shrinking. The nostrils are separated by the septum. The septum can sometimes be deviated, causing one nostril to appear larger than the other. With extreme damage to the septum and columella, the two nostrils are no longer separated and form a single larger external opening. Like other tetrapods, humans have two external nostrils (anterior nares) and two additional nostrils at the back of the nasal cavity, inside the head (posterior nares, posterior nasal apertures or choanae). They also connect the nose to the throat (the nasopharynx), aiding in respiration. Though all four nostrils were on the outside of the head of the aquatic ancestors of modern tetrapods, the nostrils for outgoing water (excurrent nostrils) migrated to the inside of the mouth, as evidenced by the discovery of Kenichthys campbelli, a 395-million-year-old fossilized lobe-finned fish which shows this migration in progress. It has two nostrils between its front teeth, similar to human embryos at an early stage. If these fail to join up, the result is a cleft palate. Each external nostril contains approximately 1,000 strands of nasal hair, which function to filter foreign particles such as pollen and dust. It is possible for humans to smell different olfactory inputs in the two nostrils and experience a perceptual rivalry akin to that of binocular rivalry when there are two different inputs to the two eyes. Furthermore, scent information from the two nostrils leads to two types of neural activity with the first cycle corresponding to the ipsilateral and the second cycle corresponding to the contralateral odor representations. In some cultures the extreme wide flaring of the nostrils accompanied by the baring of the upper teeth is often referred to as "doing the nostrils." The Procellariiformes are distinguished from other birds by having tubular extensions of their nostrils. Widely-spaced nostrils, like those of the hammerhead shark, may be useful in determining the direction of an odour's source.	Biology and health sciences	Sensory nervous system	Biology
69053	https://en.wikipedia.org/wiki/Creosote	Creosote	Creosote is a category of carbonaceous chemicals formed by the distillation of various tars and pyrolysis of plant-derived material, such as wood, or fossil fuel. They are typically used as preservatives or antiseptics. Some creosote types were used historically as a treatment for components of seagoing and outdoor wood structures to prevent rot (e.g., bridgework and railroad ties, see image). Samples may be found commonly inside chimney flues, where the coal or wood burns under variable conditions, producing soot and tarry smoke. Creosotes are the principal chemicals responsible for the stability, scent, and flavor characteristic of smoked meat; the name is derived . The two main kinds recognized in industry are coal-tar creosote and wood-tar creosote. The coal-tar variety, having stronger and more toxic properties, has chiefly been used as a preservative for wood; coal-tar creosote was also formerly used as an escharotic, to burn malignant skin tissue, and in dentistry, to prevent necrosis, before its carcinogenic properties became known. The wood-tar variety has been used for meat preservation, ship treatment, and such medical purposes as an anaesthetic, antiseptic, astringent, expectorant, and laxative, though these have mostly been replaced by modern formulations. Varieties of creosote have also been made from both oil shale and petroleum, and are known as oil-tar creosote when derived from oil tar, and as water-gas-tar creosote when derived from the tar of water gas. Creosote also has been made from pre-coal formations such as lignite, yielding lignite-tar creosote, and peat, yielding peat-tar creosote. Creosote oils The term creosote has a broad range of definitions depending on the origin of the coal tar oil and end-use of the material. With respect to wood preservatives, the United States Environmental Protection Agency (EPA) considers the term creosote to mean a pesticide for use as a wood preservative meeting the American Wood Protection Association (AWPA) Standards P1/P13 and P2. The AWPA Standards require that creosote "shall be a pure coal tar product derived entirely from tar produced by the carbonization of bituminous coal." Currently, all creosote-treated wood products—foundation and marine pilings, lumber, posts, railroad ties, timbers, and utility poles—are manufactured using this type of wood preservative. The manufacturing process can only be a pressure process under the supervision of a licensed applicator certified by the State Departments of Agriculture. No brush-on, spray, or non-pressure uses of creosote are allowed, as specified by the EPA-approved label for the use of creosote. The use of creosote according to the AWPA Standards does not allow for mixing with other types of "creosote type" materials—such as lignite-tar creosote, oil-tar creosote, peat-tar creosote, water-gas-tar creosote, or wood-tar creosote. The AWPA Standard P3 does however, allow blending of a high-boiling petroleum oil meeting the AWPA Standard P4. The information that follows describing the other various types of creosote materials and its uses should be considered as primarily being of only historical value. This history is important, because it traces the origin of these different materials used during the 19th and early 20th centuries. Furthermore, it must be considered that these other types of creosotes – lignite-tar, wood-tar, water-gas-tar, etc. – are not currently being manufactured and have either been replaced with more-economical materials, or replaced by products that are more efficacious or safer. For some part of their history, coal-tar creosote and wood-tar creosote were thought to have been equivalent substances—albeit of distinct origins—accounting for their common name; the two were determined only later to be chemically different. All types of creosote are composed of phenol derivatives and share some quantity of monosubstituted phenols, but these are not the only active element of creosote. For their useful effects, coal-tar creosote relies on the presence of naphthalenes and anthracenes, while wood-tar creosote relies on the presence of methyl ethers of phenol. Otherwise, either type of tar would dissolve in water. Creosote was first discovered in its wood-tar form in 1832, by Carl Reichenbach, when he found it both in the tar and in pyroligneous acids obtained by a dry distillation of beechwood. Because pyroligneous acid was known as an antiseptic and meat preservative, Reichenbach conducted experiments by dipping meat in a dilute solution of distilled creosote. He found that the meat was dried without undergoing putrefaction and had attained a smoky flavor. This led him to reason that creosote was the antiseptic component contained in smoke, and he further argued that the creosote he had found in wood tar was also in coal tar, as well as amber tar and animal tar, in the same abundance as in wood tar. Soon afterward, in 1834, Friedrich Ferdinand Runge discovered carbolic acid (phenol) in coal-tar, and Auguste Laurent obtained it from "phenylhydrate", which was soon determined to be the same compound. There was no clear view on the relationship between carbolic acid and creosote; Runge described it as having similar caustic and antiseptic properties, but noted that it was different, in that it was an acid and formed salts. Nonetheless, Reichenbach argued that creosote was also the active element, as it was in pyroligneous acid. Despite evidence to the contrary, his view held sway with most chemists, and it became commonly accepted wisdom that creosote, carbolic acid, and phenylhydrate were identical substances, with different degrees of purity. Carbolic acid was soon commonly sold under the name "creosote", and the scarcity of wood-tar creosote in some places led chemists to believe that it was the same substance as that described by Reichenbach. In the 1840s, Eugen Freiherr von Gorup-Besanez, after realizing that two samples of substances labelled as creosote were different, started a series of investigations to determine the chemical nature of carbolic acid, leading to a conclusion that it more resembled chlorinated quinones and must have been a different, entirely unrelated substance. Independently, there were investigations into the chemical nature of creosote. A study by F.K. Völkel revealed that the smell of purified creosote resembled that of guaiacol, and later studies by Heinrich Hlasiwetz identified a substance common to guaiacum and creosote that he called creosol, and he determined that creosote contained a mixture of creosol and guaiacol. Later investigations by Gorup-Besanez, A.E. Hoffmann, and Siegfried Marasse showed that wood-tar creosote also contained phenols, giving it a feature in common with coal-tar creosote. Historically, coal-tar creosote has been distinguished from what was thought of as creosote proper—the original substance of Reichenbach's discovery—and it has been referred to specifically as "creosote oil". But, because creosote from coal-tar and wood-tar are obtained from a similar process and have some common uses, they have also been placed in the same class of substances, with the terms "creosote" or "creosote oil" referring to either product. Wood-tar creosote Wood-tar creosote is a colourless to yellowish greasy liquid with a smoky odor, produces a sooty flame when burned, and has a burned taste. It is non-buoyant in water, with a specific gravity of 1.037 to 1.087, retains fluidity at a very low temperature, and boils at 205-225 °C. In its purest form, it is transparent. Dissolution in water requires up to 200 times the amount of water as the base creosote. This creosote is a combination of natural phenols: primarily guaiacol and creosol (4-methylguaiacol), which typically constitutes 50% of the oil; second in prevalence are cresol and xylenol; the rest being a combination of monophenols and polyphenols. The simple phenols are not the only active element in wood-tar creosote. In solution, they coagulate albumin, which is a water-soluble protein found in meat, so they serve as a preserving agent, but also cause denaturation. Most of the phenols in the creosote are methoxy derivatives: they contain the methoxy group linked to the benzene nucleus. The high level of methyl derivates created from the action of heat on wood (also apparent in the methyl alcohol produced through distillation) make wood-tar creosote substantially different from coal-tar creosote. Guaiacol is a methyl ether of pyrocatechin, while creosol is a methyl ether of methyl-pyrocatechin, the next homolog of pyrocatechin. Methyl ethers differ from simple phenols in being less hydrophilic, caustic, and poisonous. This allows meat to be successfully preserved without tissue denaturation, and allows creosote to be used as a medical ointment. Because wood-tar creosote is used for its guaiacol and creosol content, it is generally derived from beechwood rather than other woods, since it distills with a higher proportion of those chemicals to other phenolics. The creosote can be obtained by distilling the wood tar and treating the fraction heavier than water with a sodium hydroxide solution. The alkaline solution is then separated from the insoluble oily layer, boiled in contact with air to reduce impurities, and decomposed by diluted sulfuric acid. This produces a crude creosote, which is purified by re-solution in alkali, re-precipitation with acid, then redistilled with the fraction passing over between 200° and 225° constituting the purified creosote. When ferric chloride is added to a dilute solution, it will turn green: a characteristic of ortho-oxy derivatives of benzene. It dissolves in sulfuric acid to a red liquid, which slowly changes to purple-violet. Shaken with hydrochloric acid in the absence of air, it becomes red, the color changing in the presence of air to dark brown or black. In preparation of food by smoking, guaiacol contributes mainly to the smoky taste, while the dimethyl ether of pyrogallol, syringol, is the main chemical responsible for the smoky aroma. Historical uses Industrial Soon after it was discovered and recognized as the principle of meat smoking, wood-tar creosote became used as a replacement for the process. Several methods were used to apply the creosote. One was to dip the meat in pyroligneous acid or a water of diluted creosote, as Reichenbach did, or brush it over with them, and within one hour the meat would have the same quality of that of traditionally smoked preparations. Sometimes the creosote was diluted in vinegar rather than water, as vinegar was also used as a preservative. Another was to place the meat in a closed box, and place with it a few drops of creosote in a small bottle. Because of the volatility of the creosote, the atmosphere was filled with a vapour containing it, and it would cover the flesh. The application of wood tar to seagoing vessels was practiced through the 18th century and early 19th century, before the creosote was isolated as a compound. Wood-tar creosote was found not to be as effective in wood treatments, because it was harder to infuse the creosote into the wood cells, but still experiments were done, including by many governments, because it proved to be less expensive on the market. Medical Even before creosote as a chemical compound was discovered, it was the chief active component of medicinal remedies in different cultures around the world. In antiquity, pitches and resins were used commonly as medicines. Pliny mentions a variety of tar-like substances being used as medicine, including cedria and pissinum. Cedria was the pitch and resin of the cedar tree, being equivalent to the oil of tar and pyroligneous acid which are used in the first stage of distilling creosote. He recommends cedria to ease the pain in a toothache, as an injection in the ear in case of hardness of hearing, to kill parasitic worms, as a preventive for infusion, as a treatment for phthiriasis and porrigo, as an antidote for the poison of the sea hare, as a liniment for elephantiasis, and as an ointment to treat ulcers both on the skin and in the lungs. He further speaks of cedria being used as the embalming agent for preparing mummies. Pissinum was a tar water that was made by boiling cedria, spreading wool fleeces over the vessels to catch the steam, and then wringing them out. The Pharmacopée de Lyon, published in 1778, says that cedar tree oil is believed to cure vomiting and help medicate tumors and ulcers. Physicians contemporary to the discovery of creosote recommended ointments and pills made from tar or pitch to treat skin diseases. Tar water had been used as a folk remedy since the Middle Ages to treat affections like dyspepsia. Bishop Berkeley wrote several works on the medical virtues of tar water, including a philosophical work in 1744 titled Siris: a chain of philosophical reflexions and inquiries concerning the virtues of tar water, and divers other subjects connected together and arising one from another, and a poem where he praised its virtues. Pyroligneous acid was also used at the time in a medicinal water called Aqua Binelli (Binelli's water), a compound which its inventor, the Italian Fedele Binelli, claimed to have hemostatic properties in his research published in 1797. These claims have since been disproven. Given this history, and the antiseptic properties known to creosote, it became popular among physicians in the 19th century. A dilution of creosote in water was sold in pharmacies as Aqua creosoti, as suggested by the previous use of pyroligneous acid. It was prescribed to quell the irritability of the stomach and bowels and detoxify, treat ulcers and abscesses, neutralize bad odors, and stimulate the mucous tissues of the mouth and throat. Creosote in general was listed as an irritant, styptic, antiseptic, narcotic, and diuretic, and in small doses when taken internally as a sedative and anaesthetic. It was used to treat ulcers, and as a way to sterilize the tooth and deaden the pain in case of a tooth-ache. Creosote was suggested as a treatment for tuberculosis by Reichenbach as early as 1833. Following Reichenbach, it was argued for by John Elliotson and Sir John Rose Cormack. Elliotson, inspired by the use of creosote to arrest vomiting during an outbreak of cholera, suggested its use for tuberculosis through inhalation. He also suggested it for epilepsy, neuralgia, diabetes, and chronic glanders. The idea of using it for tuberculosis failed to be accepted. Use for this purpose was dropped, until the idea was revived in 1876 by British doctor G. Anderson Imlay, who suggested it be applied locally by spray to the bronchial mucous membrane. This was followed up in 1877 when it was argued for in a clinical paper by Charles Bouchard and Henri Gimbert. Germ theory had been established by Pasteur in 1860, and Bouchard, arguing that a bacillus was responsible for the disease, sought to rehabilitate creosote for its use as an antiseptic to treat it. He began a series of trials with Gimbert to convince the scientific community, and claimed a promising cure rate. A number of publications in Germany confirmed his results in the following years. Later, a period of experimentation with different techniques and chemicals using creosote in treating tuberculosis lasted until about 1910, when radiation therapy seemed more promising. Guaiacol, instead of a full creosote solution, was suggested by Hermann Sahli in 1887. He argued it had the active chemical of creosote and had the advantage of being of definite composition and having a less unpleasant taste and odor. A number of solutions of both creosote and guaiacol appeared on the market, such as phosphotal and guaicophosphal, phosphites of creosote and guaiacol; eosot and geosot, valerinates of creosote and guaicol; phosot and taphosot, phosphate and tannophospate of creosote; and creosotal and tanosal, tannates of creosote. Creosote and eucalyptus oil were also a remedy used together, administered through a vaporizor and inhaler. Since then, more effective and safer treatments for tuberculosis have been developed. In the 1940s, Canadian-based Eldon Boyd experimented with guaiacol and a recent synthetic modification—glycerol guaiacolate (guaifenesin)—on animals. His data showed that both drugs were effective in increasing secretions into the airways in laboratory animals, when high-enough doses were given. Current uses Industrial Wood-tar creosote is to some extent used for wood preservation, but it is generally mixed with coal-tar creosote, since the former is not as effective. Commercially available preparations of "liquid smoke", marketed to add a smoked flavour to meat and aid as a preservative, consist primarily of creosote and other constituents of smoke. Creosote is the ingredient that gives liquid smoke its function; guaicol lends to the taste and the creosote oils help act as the preservative. Creosote can be destroyed by treatment with chlorine, either sodium hypochlorite, or calcium hypochlorite solutions. The phenol ring is essentially opened, and the molecule is then subject to normal digestion and normal respiration. Medical The guaifenesin developed by Eldon Boyd is still commonly used today as an expectorant, sold over the counter, and usually taken by mouth to assist the bringing up of phlegm from the airways in acute respiratory tract infections. Guaifenesin is a component of Mucinex, Robitussin DAC, Cheratussin DAC, Robitussin AC, Cheratussin AC, Benylin, DayQuil Mucous Control, Meltus, and Bidex 400. Seirogan is a popular Kampo medicine in Japan, used as an anti-diarrheal, and has 133 mg wood creosote from beech, pine, maple or oak wood per adult dose as its primary ingredient. Seirogan was first used as a gastrointestinal medication by the Imperial Japanese Army in Russia during the Russo-Japanese War of 1904 to 1905. Creomulsion is a cough medicine in the United States, introduced in 1925, that is still sold and contains beechwood creosote. Beechwood creosote is also found under the name kreosotum or kreosote. Coal-tar creosote Coal-tar creosote is greenish-brown liquid, with different degrees of darkness, viscosity, and fluorescence depending on how it is made. When freshly made, the creosote is a yellow oil with a greenish cast and highly fluorescent, and the fluorescence is increased by exposure to air and light. After settling, the oil is dark green by reflected light and dark red by transmitted light. To the naked eye, it generally appears brown. The creosote (often called "creosote oil") consists almost wholly of aromatic hydrocarbons, with some amount of bases and acids and other neutral oils. The flash point is 70–75 °C and burning point is 90–100 °C, and when burned it releases a greenish smoke. The smell largely depends on the naphtha content in the creosote. If there is a high amount, it will have a naphtha-like smell, otherwise it will smell more of tar. In the process of coal-tar distillation, the distillate is collected into four fractions; the "light oil", which remains lighter than water, the "middle oil" which passes over when the light oil is removed; the "heavy oil", which sinks; and the "anthracene oil", which when cold is mostly solid and greasy, of a buttery consistence. Creosote refers to the portion of coal tar which distills as "heavy oil", typically between 230 and 270 °C, also called "dead oil"; it sinks into water but still is fairly liquid. Carbolic acid is produced in the second fraction of distillation and is often distilled into what is referred to as "carbolic oil". Commercial creosote contains substances from six groups. The two groups occur in the greatest amounts and are the products of the distillation process—the "tar acids", which distill below 205 °C and consist mainly of phenols, cresols, and xylenols, including carbolic acid—and aromatic hydrocarbons, which divide into naphthalenes, which distill approximately between 205 and 255 °C, and constituents of an anthracene nature, which distill above 255 °C. The quantity of each varies based on the quality of tar and temperatures used, but generally, the tar acids won't exceed 5%, the naphthalenes make up 15 to 50%, and the anthracenes make up 45% to 70%. The hydrocarbons are mainly aromatic; derivatives of benzene and related cyclic compounds such as naphthalene, anthracene, phenanthrene, acenaphthene, and fluorene. Creosotes from vertical-retort and low temperature tars contain, in addition, some paraffinic and olefinic hydrocarbons. The tar-acid content also depends on the source of the tar—it may be less than 3% in creosote from coke-oven tar and as high as 32% in creosote from vertical retort tar. All of these have antiseptic properties. The tar acids are the strongest antiseptics but have the highest degree of solubility in water and are the most volatile; so, like with wood-tar creosote, phenols are not the most valued component, as by themselves they would lend to being poor preservatives. In addition, creosote contains several products naturally occurring in coal—nitrogen-containing heterocycles, such as acridines, carbazoles, and quinolines, referred to as the "tar bases" and generally make up about 3% of the creosote—sulfur-containing heterocycles, generally benzothiophenes—and oxygen-containing heterocycles, dibenzofurans. Lastly, creosote contains a small number of aromatic amines produced by the other substances during the distillation process and likely resulting from a combination of thermolysis and hydrogenation. The tar bases are often extracted by washing the creosote with aqueous mineral acid, although they're also suggested to have antiseptic ability similar to the tar acids. Commercially used creosote is often treated to extract the carbolic acid, naphthalene, or anthracene content. The carbolic acid or naphthalene is generally extracted to be used in other commercial products. In the early 20th century, American-produced creosote oils typically had low amounts of anthracene and high amounts of naphthalene, because when forcing the distillate at a temperature that produces anthracene the soft pitch will be ruined and only the hard pitch will remain; this ruined it for use in roofing purposes (which was common before widespread availability of cheap oil bitumen) and only left a product which wasn't commercially useful. Historical uses Industrial The use of coal-tar creosote on a commercial scale began in 1838, when a patent covering the use of creosote oil to treat timber was taken out by inventor John Bethell. The "Bethell process"—or as it later became known, the full-cell process—involves placing wood to be treated in a sealed chamber and applying a vacuum to remove air and moisture from wood "cells". The wood is then pressure-treated to imbue it with creosote or other preservative chemicals, after which vacuum is reapplied to separate the excess treatment chemicals from the timber. Alongside the zinc chloride-based "Burnett process", use of creosoted wood prepared by the Bethell process became a principal way of preserving railway timbers (most notably railway sleepers) to increase the lifespan of the timbers, and avoiding having to regularly replace them. Besides treating wood, it was also used for lighting and fuel. In the beginning, it was only used for lighting needed in harbour and outdoor work, where the smoke that was produced from burning it was of little inconvenience. By 1879, lamps had been created that ensured a more complete combustion by using compressed air, removing the drawback of the smoke. Creosote was also processed into gas and used for lighting that way. As a fuel, it was used to power ships at sea and blast furnaces for different industrial needs, once it was discovered to be more efficient than unrefined coal or wood. It was also used industrially for the softening of hard pitch, and burned to produce lamp black. By 1890, the production of creosote in the United Kingdom totaled approximately 29,900,000 gallons per year. In 1854, Alexander McDougall and Robert Angus Smith developed and patented a product called McDougall's Powder as a sewer deodorant; it mainly consisted of carbolic acid derived from creosote. McDougall, in 1864, experimented with his solution to remove entozoa parasites from cattle pasturing on a sewage farm. This later led to widespread use of creosote as a cattle wash and sheep dip. External parasites would be killed in a creosote diluted dip, and drenching tubes would be used to administer doses to the animals' stomachs to kill internal parasites. Creosoted wood blocks were a common road-paving material in the late 19th and early 20th centuries, but ultimately fell out of favor because they did not generally hold up well enough over time. Two later methods for creosoting wood were introduced after the turn of the century, referred to as empty-cell processes, because they involve compressing the air inside the wood so that the preservative can only coat the inner cell walls rather than saturating the interior cell voids. This is a less effective, though usually satisfactory, method of treating the wood, but is used because it requires less of the creosoting material. The first method, the "Rüping process" was patented in 1902, and the second, the "Lowry process" was patented in 1906. Later in 1906, the "Allardyce process" and "Card process" were patented to treat wood with a combination of both creosote and zinc chloride. In 1912, it was estimated that a total of 150,000,000 gallons were produced in the US per year. Medical Coal-tar creosote, despite its toxicity, was used as a stimulant and escharotic, as a caustic agent used to treat ulcers and malignancies, cauterize wounds, and prevent infection and decay. It was particularly used in dentistry to destroy tissues and arrest necrosis. Current uses Industrial Coal-tar creosote is the most widely used wood treatment today; both industrially, processed into wood using pressure methods such as "full-cell process" or "empty-cell process", and more commonly applied to wood through brushing. In addition to toxicity to fungi, insects, and marine borers, it serves as a natural water repellent. It is commonly used to preserve and waterproof railroad ties, pilings, telephone poles, power line poles, marine pilings, and fence posts. Although suitable for use in preserving the structural timbers of buildings, it is not generally used that way because it is difficult to apply. There are also concerns about the environmental impact of the leaching of creosote into aquatic ecosystems. Due to its carcinogenic character, the European Union has regulated the quality of creosote for the EU market and requires that the sale of creosote be limited to professional users. The United States Environmental Protection Agency regulates the use of coal-tar creosote as a wood preservative under the provisions of the Federal Insecticide, Fungicide, and Rodenticide Act. Creosote is considered a restricted-use pesticide and is only available to licensed pesticide applicators. Oil-tar creosote Oil-tar creosote is derived from the tar that forms when using petroleum or shale oil in the manufacturing of gas. The distillation of the tar from the oil occurs at very high temperatures; around 980 °C. The tar forms at the same time as the gas, and once processed for creosotes contains a high percentage of cyclic hydrocarbons, a very low amount of tar acids and tar bases, and no true anthracenes have been identified. Historically, this has mainly been produced in the United States on the Pacific coast, where petroleum has been more abundant than coal. Limited quantities have been used industrially, either alone, mixed with coal-tar creosote, or fortified with pentachlorophenol. Water-gas-tar creosote Water-gas-tar creosote is also derived from petroleum oil or shale oil, but by a different process; it is distilled during the production of water gas. The tar is a by-product resulting from enrichment of water gas with gases produced by thermal decomposition of petroleum. Of the creosotes derived from oil, it is practically the only one used for wood preservation. It has the same degree of solubility as coal-tar creosote and is easy to infuse into wood. Like standard oil-tar creosote, it has a low amount of tar acids and tar bases, and has less antiseptic qualities. Petri dish tests have shown that water-gas-tar creosote is one-sixth as anti-septically effective as that of coal-tar. Lignite-tar creosote Lignite-tar creosote is produced from lignite rather than bituminous coal, and varies considerably from coal-tar creosote. Also called "lignite oil", it has a very high content of tar acids, and has been used to increase the tar acids in normal creosote when necessary. When it has been produced, it has generally been applied in mixtures with coal-tar creosote or petroleum. Its effectiveness when used alone has not been established. In an experiment with southern yellow pine fence posts in Mississippi, straight lignite-tar creosote was giving good results after about 27 years exposure, although not as good as the standard coal-tar creosote used in the same situation. Peat-tar creosote There have also been attempts to distill creosote from peat-tar, although mostly unsuccessful due to the problems with winning and drying peat on an industrial scale. Peat tar by itself has in the past been used as a wood preservative. Health effects According to the Agency for Toxic Substances and Disease Registry (ATSDR), eating food or drinking water contaminated with high levels of coal-tar creosote may cause a burning in the mouth and throat, and stomach pains. ATSDR also states that brief direct contact with large amounts of coal-tar creosote may result in a rash or severe irritation of the skin, chemical burns of the surfaces of the eyes, convulsions and mental confusion, kidney or liver problems, unconsciousness, and even death. Longer direct skin contact with low levels of creosote mixtures or their vapours can result in increased light sensitivity, damage to the cornea, and skin damage. Longer exposure to creosote vapours can cause irritation of the respiratory tract. The International Agency for Research on Cancer (IARC) has determined that coal-tar creosote is probably carcinogenic to humans, based on adequate animal evidence and limited human evidence. The animal testing relied upon by IARC involved the continuous application of creosote to the shaved skin of rodents. After weeks of creosote application, the animals developed cancerous skin lesions and in one test, lesions of the lung. The United States Environmental Protection Agency has stated that coal-tar creosote is a probable human carcinogen based on both human and animal studies. As a result, the Federal Occupational Safety and Health Administration (OSHA) has set a permissible exposure limit of 0.2 milligrams of coal-tar creosote per cubic meter of air (0.2 mg/m3) in the workplace during an 8-hour day, and the Environmental Protection Agency (EPA) requires that spills or accidental releases into the environment of one pound (0.454 kg) or more of creosote be reported to them. There is no unique exposure pathway of children to creosote. Children exposed to creosote probably experience the same health effects seen in adults exposed to creosote. It is unknown whether children differ from adults in their susceptibility to health effects from creosote. A 2005 mortality study of creosote workers found no evidence supporting an increased risk of cancer death, as a result of exposure to creosote. Based on the findings of the largest mortality study to date of workers employed in creosote wood treating plants, there is no evidence that employment at creosote wood-treating plants or exposure to creosote-based preservatives was associated with any significant mortality increase from either site-specific cancers or non-malignant diseases. The study consisted of 2,179 employees at eleven plants in the United States where wood was treated with creosote preservatives. Some workers began work in the 1940s to 1950s. The observation period of the study covered 1979–2001. The average length of employment was 12.5 years. One third of the study subjects were employed for over 15 years. The largest health effect of creosote is deaths caused by residential chimney fires due to chimney tar (creosote) build-up. This is entirely unconnected with its industrial production or use. Build-up in chimneys Burning wood and fossil fuels in the absence of adequate airflow (such as in an enclosed furnace or stove), causes incomplete combustion of the oils in the wood, which are off-gassed as volatiles in the smoke. As the smoke rises through the chimney it cools, causing water, carbon, and volatiles to condense on the interior surfaces of the chimney flue. The black oily residue that builds up is referred to as creosote, which is similar in composition to the commercial products by the same name, but with a higher content of carbon black. Over the course of a season creosote deposits can become several inches thick. This creates a compounding problem, because the creosote deposits reduce the draft (airflow through the chimney) which increases the probability that the wood fire is not getting enough air for complete combustion. Since creosote is highly combustible, a thick accumulation creates a fire hazard. If a hot fire is built in the stove or fireplace, and the air control left wide open, this may allow hot oxygen into the chimney where it comes in contact with the creosote which then ignites—causing a chimney fire. Chimney fires often spread to the main building because the chimney gets so hot that it ignites any combustible material in direct contact with it, such as wood. The fire can also spread to the main building from sparks emitting from the chimney and landing on combustible roof surfaces. In order to properly maintain chimneys and heaters that burn wood or carbon-based fuels, the creosote buildup must be removed. Chimney sweeps perform this service for a fee. Release into environment Even though creosote is pressurized into the wood, the release of the chemical – and resulting marine pollution – occurs due to many different events: During the lifetime of the marine piling, weathering occurs from tides and water flow which slowly opens the oily outer coating and exposes the smaller internal pores to more water flow. Frequent weathering occurs daily, but more severe weather, such as hurricanes, can cause damage or loosening of the wooden pilings. Many pilings are either broken into pieces from debris, or are completely washed away during these storms. When the pilings are washed away, they come to settle on the bottom of the body of water where they reside, and then they leach chemicals into the water slowly over a long period of time. This long-term secretion is not normally noticed because the piling is submerged beneath the surface, hidden from sight. The creosote is mostly insoluble in water, but the lower-molecular-weight compounds will become soluble the longer the broken wood is exposed to the water. In this case, some of the chemicals become water-soluble and further leach into the aquatic sediment while the rest of the insoluble chemicals remain together in a tar-like substance. Another source of damage comes from wood-boring fauna, such as shipworms and Limnoria. Though creosote is used as a pesticide preservative, studies have shown that Limnoria is resistant to wood preservative pesticides and can cause small holes in the wood, through which creosote can then be released. Chemical reactions with sediment and organisms Once the soluble compounds from the creosote preservative leach into the water, the compounds begin reacting with the external environment or are consumed by organisms. The reactions vary depending on the concentration of each compound that is released from the creosote, but major reactions are outlined below: Alkylation Alkylation occurs when a molecule replaces a hydrogen atom with an alkyl group that generally comes from an organic molecule. Alkyl groups that are found naturally occurring in the environment are organometallic compounds. Organometallic compounds generally contain a methyl, ethyl, or butyl derivative which is the alkyl group that replaces the hydrogen. Other organic compounds, such as methanol, can provide alkyl groups for alkylation. Methanol is found naturally in the environment in small concentrations, and has been linked to the release from biological decomposition of waste and even a byproduct of vegetation. The following reactions are alkylations of soluble compounds found in creosote preservatives with methanol. m-Cresol The diagram above depicts a reaction between m-cresol and methanol where a c-alkylation product is produced. The c-alkylation reaction means that instead of replacing the hydrogen atom on the -OH group, the methyl group (from the methanol) replaces the hydrogen on a carbon in the benzene ring. The products of this c-alkylation can be in either a para- or ortho- orientation on the molecule, as seen in the diagram, and water, which is not shown. Isomers of the dimethylphenol (DMP) compound are the products of the para- and ortho-c-alkylation. Dimethylphenol (DMP) compound is listed as an aquatic hazard by characteristic, and is toxic with long lasting effects. Phenol This diagram shows an o-alkylation between phenol and methanol. Unlike the c-alkylation, the o-alkylation replaces the hydrogen atom on the -OH group with the methyl group (from the methanol). The product of the o-alkylation is methoxybenzene, better-known as anisole, and water, which is not shown in the diagram. Anisole is listed as an acute hazard to aquatic life with long-term effects. Bioaccumulation Bioaccumulation is the process by which an organism takes in chemicals through ingestion, exposure, and inhalation. Bioaccumulation is broken down into bioconcentration (uptake of chemicals from the environment) and biomagnification (increasing concentration of chemicals as they move up the food chain). Certain species of aquatic organisms are affected differently from the chemicals released from creosote preservatives. One of the more studied organisms is a mollusk. Mollusks attach to the wooden, marine pilings and are in direct contact with the creosote preservatives. Many studies have been conducted using polycyclic aromatic hydrocarbons (PAH), which are low molecular hydrocarbons found in some creosote-based preservatives. In a study conducted from Pensacola, Florida, a group of native mollusks were kept in a controlled environment, and a different group of native mollusks were kept in an environment contaminated with creosote preservatives. The mollusks in the contaminated environment were shown to have a bioaccumulation of up to ten times the concentration of PAH than the control species. The intake of organisms is dependent on whether the compound is in an ionized or an un-ionized form. To determine whether the compound is ionized or un-ionized, the pH of the surrounding environment must be compared to the pKa or acidity constant of the compound. If the pH of the environment is lower than the pKa, then the compound is un-ionized which means that the compound will behave as if it is non-polar. Bioaccumulation for un-ionized compounds comes from partitioning equilibrium between the aqueous phase and the lipids in the organism. If the pH is higher than the pKa, then the compound is considered to be in the ionized form. The un-ionized form is favored because the bioaccumulation is easier for the organism to intake through partitioning equilibrium. The table below shows a list of pKas from compounds found in creosote preservatives and compares them to the average pH of seawater (reported to be 8.1). Each of the compounds in the table above is found in creosote preservatives; all are in the favored un-ionized form. In another study, various species of small fish were tested to see how the exposure time to PAH chemicals affected the fish. This study showed that an exposure time of 24–96 hours on various shrimp and fish species affected the growth, reproduction, and survival functions of the organisms for most of the compounds tested. Biodegradation It can be seen in some studies that biodegradation accounts for the absence of creosote preservatives on the initial surface of the sediment. In a study from Pensacola, Florida, PAHs were not detected on the surface on the aquatic sediment, but the highest concentrations were detected at a depth of 8-13 centimeters. A form an anaerobic biodegradation of m-cresol was seen in a study using sulfate-reducing and nitrate-reducing enriched environments. The reduction of m-cresol in this study was seen in under 144 hours, while additional chemical intermediates were being formed. The chemical intermediates were formed in the presence of bicarbonate. The products included 4-hydroxy-2-methylbenzoic acid and acetate compounds. Although the conditions were enriched with the reducing anaerobic compounds, sulfate and nitrate reducing bacteria are commonly found in the environment. For further information, see sulfate-reducing bacteria. The type of anaerobic bacteria ultimately determines the reduction of the creosote preservative compounds, while each individual compound may only go through reduction under certain conditions. BTEX is a mixture of benzene, toluene, ethylbenzene, and xylene, that was studied in the presence of four different anaerobic-enriched sediments. Though the compound, BTEX, is not found in creosote preservatives, the products of creosote preservatives' oxidation-reduction reactions include some of these compounds. For oxidation-reduction reactions, see the following section. In this study, it was seen that certain compounds such as benzene were only reduced under sulfate-enriched environments, while toluene was reduced under a variety of bacteria-enriched environments, not just sulfate. The biodegradation of a creosote preservative in an anaerobic enrichment depends not only on the type of bacteria enriching the environment, but also the compound that has been released from the preservative. In aerobic environments, preservative compounds are limited in the biodegradation process by the presence of free oxygen. In an aerobic environment, free oxygen comes from oxygen saturated sediments, sources of precipitation, and plume edges. The free oxygen allows for the compounds to be oxidized and decomposed into new intermediate compounds. Studies have shown that when BTEX and PAH compounds were placed in aerobic environments, the oxidation of the ring structures caused cleavage in the aromatic ring and allowed for other functional groups to attach. When an aromatic hydrocarbon was introduced to the molecular oxygen in experimental conditions, a dihydrodiol intermediate was formed, and then oxidation occurred transforming the aromatic into a catechol compound. Catechol allows for cleavage of the aromatic ring to occur, where functional groups can then add in an ortho- or meta- position. Oxidation-reduction Even though many studies conduct testing under experimental or enriched conditions, oxidation-reduction reactions occur naturally and allow for chemicals to go through processes such as biodegradation, outlined above. Oxidation is defined as the loss of an electron to another species, while reduction is the gaining of an electron from another species. As compounds go through oxidation and reduction in sediments, the preservative compounds are altered to form new chemicals, leading to decomposition. An example of the oxidation of p-cresol and phenol can be seen in the figures below: p-Cresol This reaction shows the oxidation of p-cresol in a sulfate-enriched environment. P-cresol was seen to be the easiest to degrade through the sulfate-enriched environment, while m-cresol and o-cresol where inhibited. In the chart above, p-cresol was oxidized under an anaerobic sulfate reducing condition and formed four different intermediates. After the formation of the intermediates, the study reported further degradation of the intermediates leading to the production of carbon dioxide and methane. The p-hydroxylbenzyl alcohol, p-hydroxylbenzaldehye, p-hyrdoxylbenzoate, and benzoate intermediates all are produced from this oxidation and released into the sediments. Similar results were also produced by different studies using other forms of oxidation such as: iron-reducing organisms, Copper/Manganese Oxide catalyst, and nitrate- reducing conditions. Phenol This reaction shows the oxidation of phenol by iron and peroxide. This combination of iron, which comes from iron oxide in the sediment, and the peroxide, commonly released by animals and plants into the environment, is known as the Fenton Reagent. This reagent is used to oxidize phenol groups by the use of a radical hydroxide group produced from the peroxide in the p-benzoquinone. This product of phenol's oxidation is now leached into the environment while other products include iron(II) and water. P-benzoquinone is listed as being a very toxic, acute environmental hazard. Environmental hazards Sediment In aquatic sediments, several reactions can transform the chemicals released by the creosote preservatives into more dangerous chemicals. Most creosote preservative compounds have hazards associated with them before they are transformed. Cresol (m-, p-, and o-), phenol, guaiacol, and xylenol (1,3,4- and 1,3,5-) all are acute aquatic hazards prior to going through chemical reactions with the sediments. Alkylation reactions allows for the compounds to transition into more toxic compounds with the addition of R-groups to the major compounds found in creosote preservatives. Compounds formed through alkylation include: 3,4-dimethylphenol, 2,3-dimethylphenol, and 2,5-dimethylphenol, which are all listed as acute environmental hazards. Biodegradation controls the rate at which the sediment holds the chemicals, and the number of reactions that are able to take place. The biodegradation process can take place under many different conditions, and vary depending on the compounds that are released. Oxidation-reduction reactions allow for the compounds to be broken down into new forms of more toxic molecules. Studies have shown oxidation-reduction reactions of creosote preservative compounds included compounds that are listed as environmental hazards, such as p-benzoquinone in the oxidation of phenol. Not only are the initial compounds in creosote hazardous to the environment, but the byproducts of the chemical reactions are environmental hazardous as well. Other From the contamination of the sediment, more of the ecosystem is affected. Organisms in the sediment are now exposed to the new chemicals. Organisms are then ingested by fish and other aquatic animals. These animals now contain concentrations of hazardous chemicals which were secreted from the creosote. Other issues with ecosystems include bioaccumulation. Bioaccumulation occurs when high levels of chemicals are passed to aquatic life near the creosote pilings. Mollusks and other smaller crustaceans are at higher risk because they are directly attached to the surface of wood pilings that are filled with creosote preservative. Studies show that mollusks in these environments take on high concentrations of chemical compounds which will then be transferred through the ecosystem's food chain. Bioaccumulation contributes to the higher concentrations of chemicals within the organisms in the aquatic ecosystems.	Physical sciences	Hydrocarbons	Chemistry
69079	https://en.wikipedia.org/wiki/Ammonium	Ammonium	Ammonium is a modified form of ammonia that has an extra hydrogen atom. It is a positively charged (cationic) molecular ion with the chemical formula or . It is formed by the addition of a proton (a hydrogen nucleus) to ammonia (). Ammonium is also a general name for positively charged (protonated) substituted amines and quaternary ammonium cations (), where one or more hydrogen atoms are replaced by organic or other groups (indicated by R). Not only is ammonium a source of nitrogen and a key metabolite for many living organisms, but it is an integral part of the global nitrogen cycle. As such, human impact in recent years could have an effect on the biological communities that depend on it. Acid–base properties The ammonium ion is generated when ammonia, a weak base, reacts with Brønsted acids (proton donors): The ammonium ion is mildly acidic, reacting with Brønsted bases to return to the uncharged ammonia molecule: Thus, the treatment of concentrated solutions of ammonium salts with a strong base gives ammonia. When ammonia is dissolved in water, a tiny amount of it converts to ammonium ions: The degree to which ammonia forms the ammonium ion depends on the pH of the solution. If the pH is low, the equilibrium shifts to the right: more ammonia molecules are converted into ammonium ions. If the pH is high (the concentration of hydrogen ions is low and hydroxide ions is high), the equilibrium shifts to the left: the hydroxide ion abstracts a proton from the ammonium ion, generating ammonia. Formation of ammonium compounds can also occur in the vapor phase; for example, when ammonia vapor comes in contact with hydrogen chloride vapor, a white cloud of ammonium chloride forms, which eventually settles out as a solid in a thin white layer on surfaces. Salts and characteristic reactions Ammonium cation is found in a variety of salts such as ammonium carbonate, ammonium chloride, and ammonium nitrate. Most simple ammonium salts are very soluble in water. An exception is ammonium hexachloroplatinate, the formation of which was once used as a test for ammonium. The ammonium salts of nitrate and especially perchlorate are highly explosive, in these cases, ammonium is the reducing agent. In an unusual process, ammonium ions form an amalgam. Such species are prepared by the addition of sodium amalgam to a solution of ammonium chloride. This amalgam eventually decomposes to release ammonia and hydrogen. To find whether the ammonium ion is present in the salt, first, the salt is heated in presence of alkali hydroxide releasing a gas with a characteristic smell, which is ammonia. To further confirm ammonia, it is passed through a glass rod dipped in an solution (hydrochloric acid), creating white dense fumes of ammonium chloride. Ammonia, when passed through (copper(II) sulfate) solution, changes its color from blue to deep blue, forming Schweizer's reagent. Ammonia or ammonium ion when added to Nessler's reagent gives a brown color precipitate known as the iodide of Million's base in basic medium. Ammonium ion when added to chloroplatinic acid gives a yellow precipitate of ammonium hexachloroplatinate(IV). Ammonium ion when added to sodium cobaltinitrite gives a yellow precipitate of ammonium cobaltinitrite. Ammonium ion gives a white precipitate of ammonium bitartrate when added to potassium bitartrate. Structure and bonding The lone electron pair on the nitrogen atom (N) in ammonia, represented as a line above the N, forms a coordinate bond with a proton (). After that, all four bonds are equivalent, being polar covalent bonds. The ion has a tetrahedral structure and is isoelectronic with methane and the borohydride anion. In terms of size, the ammonium cation (rionic = 175 pm) resembles the caesium cation (rionic = 183 pm). Organic ions The hydrogen atoms in the ammonium ion can be substituted with an alkyl group or some other organic group to form a substituted ammonium ion (IUPAC nomenclature: aminium ion). Depending on the number of organic groups, the ammonium cation is called a primary, secondary, tertiary, or quaternary. Except the quaternary ammonium cations, the organic ammonium cations are weak acids. An example of a reaction forming an ammonium ion is that between dimethylamine, , and an acid to give the dimethylammonium cation, : Quaternary ammonium cations have four organic groups attached to the nitrogen atom, they lack a hydrogen atom bonded to the nitrogen atom. These cations, such as the tetra-n-butylammonium cation, are sometimes used to replace sodium or potassium ions to increase the solubility of the associated anion in organic solvents. Primary, secondary, and tertiary ammonium salts serve the same function but are less lipophilic. They are also used as phase-transfer catalysts and surfactants. An unusual class of organic ammonium salts is derivatives of amine radical cations, such as tris(4-bromophenyl)ammoniumyl hexachloroantimonate. Biology Because nitrogen often limits net primary production due to its use in enzymes that mediate the biochemical reactions that are necessary for life, ammonium is utilized by some microbes and plants. For example, energy is released by the oxidation of ammonium in a process known as nitrification, which produces nitrate and nitrite. This process is a form of autotrophy that is common amongst Nitrosomonas, Nitrobacter, Nitrosolobus, and Nitrosospira, amongst others. The amount of ammonium in soil that is available for nitrification by microbes varies depending on environmental conditions. For example, ammonium is deposited as a waste product from some animals, although it is converted into urea in mammals, sharks, and amphibians, and into uric acid in birds, reptiles, and terrestrial snails. Its availability in soils is also influenced by mineralization, which makes more ammonium available from organic nitrogen sources, and immobilization, which sequesters ammonium into organic nitrogen sources, both of which are mitigated by biological factors. Conversely, nitrate and nitrite can be reduced to ammonium as a way for living organisms to access nitrogen for growth in a process known as assimilatory nitrate reduction. Once assimilated, it can be incorporated into proteins and DNA. Ammonium can accumulate in soils where nitrification is slow or inhibited, which is common in hypoxic soils. For example, ammonium mobilization is one of the key factors for the symbiotic association between plants and fungi, called mycorrhizae. However, plants that consistently utilize ammonium as a nitrogen source often must invest into more extensive root systems due to ammonium's limited mobility in soils compared to other nitrogen sources. Human impact Ammonium deposition from the atmosphere has increased in recent years due to volatilization from livestock waste and increased fertilizer use. Because net primary production is often limited by nitrogen, increased ammonium levels could impact the biological communities that rely on it. For example, increasing nitrogen content has been shown to increase plant growth, but aggravate soil phosphorus levels, which can impact microbial communities. Metal The ammonium cation has very similar properties to the heavier alkali metal cations and is often considered a close equivalent. Ammonium is expected to behave as a metal ( ions in a sea of electrons) at very high pressures, such as inside giant planets such as Uranus and Neptune. Under normal conditions, ammonium does not exist as a pure metal but does as an amalgam (alloy with mercury).	Physical sciences	Salts and ions: General	Chemistry
69149	https://en.wikipedia.org/wiki/MMR%20vaccine	MMR vaccine	The MMR vaccine is a vaccine against measles, mumps, and rubella (German measles), abbreviated as MMR. The first dose is generally given to children around 9 months to 15 months of age, with a second dose at 15 months to 6 years of age, with at least four weeks between the doses. After two doses, 97% of people are protected against measles, 88% against mumps, and at least 97% against rubella. The vaccine is also recommended for those who do not have evidence of immunity, those with well-controlled HIV/AIDS, and within 72 hours of exposure to measles among those who are incompletely immunized. It is given by injection. The MMR vaccine is widely used around the world. Worldwide over 500 million doses were administered between 1999 and 2004, and 575 million doses have been administered since the vaccine's introduction worldwide. Measles resulted in 2.6 million deaths per year before immunization became common. This has decreased to 122,000 deaths per year mostly in low-income countries. Through vaccination, , rates of measles in North and South America are very low. Rates of disease have been seen to increase in populations that go unvaccinated. Between 2000 and 2018, vaccination decreased measles deaths by 73%. Side effects of immunization are generally mild and resolve without any specific treatment. These may include fever, as well as pain or redness at the injection site. Severe allergic reactions occur in about one in a million people. Because it contains live viruses, the MMR vaccine is not recommended during pregnancy but may be given during breastfeeding. The vaccine is safe to give at the same time as other vaccines. Being recently immunized does not increase the risk of passing measles, mumps, or rubella on to others. There is no evidence of an association between MMR immunisation and autistic spectrum disorders. The MMR vaccine is a mixture of live weakened viruses of the three diseases. The MMR vaccine was developed by Maurice Hilleman. It was licensed for use in the US by Merck in 1971. Stand-alone measles, mumps, and rubella vaccines had been previously licensed in 1963, 1967, and 1969, respectively. Recommendations for a second dose were introduced in 1989. The MMRV vaccine, which also covers chickenpox, may be used instead. An MR vaccine, without coverage for mumps, is also occasionally used. Medical use Cochrane concluded that the "Existing evidence on the safety and effectiveness of MMR and MMRV vaccine supports current policies of mass immunisation aimed at global measles eradication to reduce morbidity and mortality associated with measles mumps rubella and varicella." The combined MMR vaccine induces immunity less painfully than three separate injections at the same time, and sooner and more efficiently than three injections given on different dates. Public Health England reports that providing a single combined vaccine as of 1988, rather than giving the option to have them also done separately, increased uptake of the vaccine. Measles Before the widespread use of a vaccine against measles, rates of disease were so high that infection was felt to be "as inevitable as death and taxes." Reported cases of measles in the United States fell from hundreds of thousands to tens of thousands per year following introduction of the vaccine in 1963. Increasing uptake of the vaccine following outbreaks in 1971, and 1977, brought this down to thousands of cases per year in the 1980s. An outbreak of almost 30,000 cases in 1990 led to a renewed push for vaccination and the addition of a second vaccine to the recommended schedule. Fewer than 200 cases have been reported in the US each year between 1997 and 2013, and the disease is no longer considered endemic there. The benefit of measles vaccination in preventing illness, disability, and death has been well documented. The first 20 years of licensed measles vaccination in the US prevented an estimated 52 million cases of the disease, 17,400 cases of intellectual disability, and 5,200 deaths. During 1999–2004, a strategy led by the World Health Organization and UNICEF led to improvements in measles vaccination coverage that averted an estimated 1.4 million measles deaths worldwide. Between 2000 and 2018, measles vaccination resulted in a 73% decrease in deaths from the disease. Measles is common in many areas of the world. Although it was declared eliminated from the US in 2000, high rates of vaccination and good communication with people who refuse vaccination are needed to prevent outbreaks and sustain the elimination of measles in the US. Of the 66 cases of measles reported in the US in 2005, slightly over half were attributable to one unvaccinated individual who acquired measles during a visit to Romania. This individual returned to a community with many unvaccinated children. The resulting outbreak infected 34 people, mostly children and virtually all unvaccinated; 9% were hospitalized, and the cost of containing the outbreak was estimated at $167,685. A major epidemic was averted due to high rates of vaccination in the surrounding communities. In 2017, an outbreak of measles occurred among the Somali-American community in Minnesota, where MMR vaccination rates had declined due to the misconception that the vaccine could cause autism. The US Centers for Disease Control and Prevention recorded 65 affected children in the outbreak by April 2017. Rubella Rubella, also known as German measles, was also very common before widespread vaccination. The major risk of rubella is during pregnancy when the baby may contract congenital rubella, which can cause significant congenital defects. Mumps Mumps is another viral disease that was once very common, especially during childhood. If mumps is acquired by a male who is past puberty, a possible complication is bilateral orchitis, which can in some cases lead to sterility. Administration The MMR vaccine is administered by a subcutaneous injection, the first dose typically at twelve months of age. The second dose may be given as early as one month after the first dose. The second dose is a dose to produce immunity in the small number of persons (2–5%) who fail to develop measles immunity after the first dose. In the US it is done before entry to kindergarten because that is a convenient time. Areas where measles is common typically recommend the first dose at nine months of age and the second dose at fifteen months of age. Safety Adverse reactions, rarely serious, may occur from each component of the MMR vaccine. Ten percent of children develop fever, malaise, and a rash 5–21 days after the first vaccination; and 3% develop joint pain lasting 18 days on average. Older women appear to be more at risk of joint pain, acute arthritis, and even (rarely) chronic arthritis. Anaphylaxis is an extremely rare but serious allergic reaction to the vaccine. One cause can be egg allergy. In 2014, the FDA approved two additional possible adverse events on the vaccination label: acute disseminated encephalomyelitis (ADEM), and transverse myelitis, with permission to also add "difficulty walking" to the package inserts. A 2012 IOM report found that the measles component of the MMR vaccine can cause measles inclusion body encephalitis in immunocompromised individuals. This report also rejected any connection between the MMR vaccine and autism. Some versions of the vaccine contain the antibiotic neomycin and therefore should not be used in people allergic to this antibiotic. The number of reports on neurological disorders is very small, other than evidence for an association between a form of the MMR vaccine containing the Urabe mumps strain and rare adverse events of aseptic meningitis, a form of viral meningitis. The UK National Health Service stopped using the Urabe mumps strain in the early 1990s due to cases of transient mild viral meningitis, and switched to a form using the Jeryl Lynn mumps strain instead. The Urabe strain remains in use in a number of countries; MMR with the Urabe strain is much cheaper to manufacture than with the Jeryl Lynn strain, and a strain with higher efficacy along with a somewhat higher rate of mild side effects may still have the advantage of reduced incidence of overall adverse events. A Cochrane review found that, compared with placebo, MMR vaccine was associated with fewer upper respiratory tract infections, more irritability, and a similar number of other adverse effects. Naturally acquired measles often occurs with immune thrombocytopenic purpura (ITP, a purpuric rash and an increased tendency to bleed that resolves within two months in children), occurring in 1 to 20,000 cases. Approximately 1 in 40,000 children are thought to acquire ITP in the six weeks following an MMR vaccination. ITP below the age of six years is generally a mild disease, rarely having long-term consequences. False claims about autism In 1998 Andrew Wakefield et al. published a fraudulent paper about twelve children, reportedly with bowel symptoms and autism or other disorders acquired soon after administration of MMR vaccine, while supporting a competing vaccine. In 2010, Wakefield's research was found by the General Medical Council to have been "dishonest", and The Lancet fully retracted the paper. Three months following The Lancet's retraction, Wakefield was struck off the UK medical register, with a statement identifying deliberate falsification in the research published in The Lancet, and was barred from practising medicine in the UK. The research was declared fraudulent in 2011 by the British Medical Journal. Since Wakefield's publication, multiple peer-reviewed studies have failed to show any association between the vaccine and autism. The US Centers for Disease Control and Prevention, the Institute of Medicine of the US National Academy of Sciences, the UK National Health Service and the Cochrane Library review have all concluded that there is no evidence of a link. Administering the vaccines in three separate doses does not reduce the chance of adverse effects, and it increases the opportunity for infection by the two diseases not immunized against first. Health experts have criticized media reporting of the MMR-autism controversy for triggering a decline in vaccination rates. Before publication of Wakefield's article, the inoculation rate for MMR in the UK was 92%; after publication, the rate dropped to below 80%. In 1998, there were 56 measles cases in the UK; by 2008, there were 1348 cases, with two confirmed deaths. In Japan, the MMR triplet is not used. Immunity is achieved by a combination vaccine for measles and rubella, followed up later with a mumps only vaccine. This has had no effect on autism rates in the country, further disproving the MMR autism hypothesis. History The component viral strains of MMR vaccine were developed by propagation in animal and human cells. For example, in the case of mumps and measles viruses, the virus strains were grown in embryonated chicken eggs. This produced strains of virus which were adapted for chicken cells and less well-suited for human cells. These strains are therefore called attenuated strains. They are sometimes referred to as neuroattenuated because these strains are less virulent to human neurons than the wild strains. The rubella component, Meruvax, was developed in 1967, through propagation using the human embryonic lung cell line WI-38 (named for the Wistar Institute) that was derived six years earlier in 1961. The term "MPR vaccine" is also used to refer to this vaccine, whereas "P" refers to parotitis which is caused by mumps. Merck MMR II is supplied freeze-dried (lyophilized) and contains live viruses. Before injection, it is reconstituted with the solvent provided. According to a review published in 2018, the GlaxoSmithKline (GSK) MMR vaccine known as Pluserix "contains the Schwarz measles virus, the Jeryl Lynn–like mumps strain, and RA27/3 rubella virus". Pluserix was introduced in Hungary in 1999. Enders' Edmonston strain has been used since 1999 in Hungary in Merck MMR II product. GSK Priorix vaccine, which uses attenuated Schwarz Measles, was introduced in Hungary in 2003. MMRV vaccine The MMRV vaccine, a combined measles, mumps, rubella, and varicella (chickenpox) vaccine, has been proposed as a replacement for the MMR vaccine to simplify the administration of the vaccines. Preliminary data indicate a rate of febrile seizures of 9 per 10,000 vaccinations with MMRV, as opposed to 4 per 10,000 for separate MMR and varicella shots; US health officials therefore, do not express a preference for use of MMRV vaccine over separate injections. In a 2012 study pediatricians and family doctors were sent a survey to gauge their awareness of the increased risk of febrile seizures (fever fits) in the MMRV. 74% of family doctors and 29% of pediatricians were unaware of the increased risk of febrile seizures. After reading an informational statement only 7% of family doctors and 20% of pediatricians would recommend the MMRV for a healthy 12- to 15-month-old child. The factor that was reported as the "most important" deciding factor in recommending the MMRV over the MMR+V was ACIP/AAFP/AAP recommendations (pediatricians, 77%; family physicians, 73%). MR vaccine This is a vaccine that covers measles and rubella but not mumps. As of 2014, it was used in a "few (unidentified) countries". Society and culture Religious concerns Some brands of the vaccine use gelatin, derived from pigs, as a stabilizer. This has caused reduced take-up among some communities, despite the fact that alternative vaccines without pig derivatives are approved and available.	Biology and health sciences	Vaccines	Health
69220	https://en.wikipedia.org/wiki/Toxicodendron	Toxicodendron	Toxicodendron is a genus of flowering plants in the sumac family, Anacardiaceae. It contains trees, shrubs and woody vines, including poison ivy, poison oak, and the lacquer tree. All members of the genus produce the skin-irritating oil urushiol, which can cause a severe allergic reaction. The generic name is derived from the Greek words τοξικός (toxikos), meaning "poison," and δένδρον (dendron), meaning "tree". The best-known members of the genus in North America are eastern poison ivy (T. radicans) and western poison oak (T. diversilobum), both ubiquitous throughout much of their respective region. The resins of certain species native to Japan, China and other Asian countries, such as lacquer tree (T. vernicifluum) and wax tree (T. succedaneum), are used to make lacquer, and, as a byproduct of lacquer manufacture, their berries are used to make japan wax. Description Plants in the genus have pinnately compound, alternate leaves and whitish or grayish drupes. They are quite variable in appearance. The leaves may have smooth, toothed, or lobed edges, and all three types of leaf edges may be present in a single plant. The plants grow as creeping vines, climbing vines, shrubs, or, in the case of lacquer tree (T. vernicifluum) and poison sumac (T. vernix), as trees. While leaves of poison ivy and poison oaks usually have three leaflets, sometimes there are five or, occasionally, even seven leaflets. Leaves of poison sumac have 7–13 leaflets, and of Lacquer Tree, 7–19 leaflets. Taxonomy It was published by Philip Miller in 1754. The lectotype species is Toxicodendron pubescens The genus is a member of the Rhus complex, and has at various times been categorized as being either its own genus or a sub-genus of Rhus. There is evidence which points to keeping Toxicodendron as a separate monophyletic genus, but researchers have stated that the Toxicodendron and Rhus groups are complex and require more study to be fully understood. The common names come from similar appearances to other species that are not closely related and to the allergic response to the urushiol. Poison oak is not an oak (Quercus, family Fagaceae), but this common name comes from the leaves' resemblance to white oak (Quercus alba) leaves, while poison ivy is not an ivy (Hedera, family Araliaceae), but has a superficially similar growth form. Technically, the plants do not contain a poison; they contain a potent allergen. Species 29 species are accepted. Toxicodendron acuminatum (synonym Rhus acuminata) – China, Bhutan, India and Nepal. Toxicodendron bimannii – Assam Toxicodendron borneense – Borneo Toxicodendron calcicola – endemic to China Toxicodendron delavayi – southwestern Sichuan and northwestern and central Yunnan in south-central China Toxicodendron diversilobum (synonym Rhus diversiloba) – Western poison oak is found throughout much of western North America, ranging from the Pacific coast into the Sierra Nevada and Cascade mountain ranges between southern British Columbia and southward into Baja California. It is extremely common in that region, where it is the predominant species of the genus. Indeed, it is California's most prevalent woody shrub. Extremely variable, it grows as a dense shrub in open sunlight, or as a climbing vine in shaded areas. It propagates by creeping rhizomes or by seed. The compound leaves are divided into three leaflets, 35–100 mm long, with scalloped, toothed, or lobed edges. The leaves may be red, yellow, green, or some combination of those colors, depending on various factors, such as the time of year. Toxicodendron fulvum – southern Yunnan and northern Thailand Toxicodendron grandiflorum – Yunnan and southwestern Sichuan in south-central China Toxicodendron griffithii – eastern Himalayas to Yunnan and southwestern Guizhou in south-central China Toxicodendron hirtellum – southern Sichuan Toxicodendron hookeri – eastern Nepal to Assam Toxicodendron khasianum – Assam and Bangladesh Toxicodendron × lobadioides (T. diversilobum × T. rydbergii) – Washington in the northwestern United States Toxicodendron nodosum – western Malesia and southwestern Sulawesi Toxicodendron oligophyllum – Fujian in southeastern China Toxicodendron orientale (synonyms Rhus orientale and R. ambigua) – Asian poison ivy is very similar to the American poison ivy, and replaces it throughout east Asia (so similar that some texts treat it as just a variety of the American species). Toxicodendron pubescens (synonym Rhus toxicarium) – Atlantic poison oak grows mostly in sandy soils in eastern parts of the United States. Growing as a shrub, its leaves are in groups of three. Leaves are typically rounded or lobed and are densely-haired. Although it is often confused with the more common poison ivy, even in the scientific literature, Atlantic poison oak has small clumps of hair on the veins on the underside of the leaves, while poison ivy does not. Toxicodendron quinquefoliolatum – Guizhou in south-central China Toxicodendron radicans (synonym Rhus radicans) – Poison ivy is extremely common in some areas of North America. In the United States, it grows in all states east of the Rockies. It also grows in Central America. Appearing as a creeping vine, a climbing vine, or a shrub, it reproduces both by creeping rootstocks and by seeds. The appearance varies. Leaves, arranged in an alternate pattern, usually in groups of three, are from 20 to 50 mm long, pointed at the tip, and can be toothed, smooth, or lobed, but never serrated. Leaves may be shiny or dull, and the color varies with the season. Vines grow almost straight up rather than wrapping around their support and can grow to 8–10 m in height. In some cases, Poison ivy may entirely engulf the supporting structure, and vines may extend outward like limbs so that it appears to be a Poison ivy "tree". Toxicodendron rhetsoides – Thailand, Laos, and Vietnam Toxicodendron rostratum – southern Yunnan Toxicodendron rydbergii (synonym Rhus rydbergii) – Western poison ivy is found in northern parts of the eastern United States. It also exists in the western United States and Canada but is much less common than poison oak. It may grow as a vine or a shrub. It was once considered a subspecies of poison ivy. It does sometimes hybridize with the climbing species. Western poison ivy is found in much of western and central United States and Canada, although not on the West Coast. In the eastern United States, it is rarely found south of New England. Toxicodendron striatum (synonym Rhus striata) – Manzanillo is a South American poisonous tree growing in the tropical rain forests on low elevation slopes. The name should not be confused with the unrelated Manchineel, a poisonous tree that is not a member of the Anacardiaceae. Toxicodendron succedaneum (synonym Rhus succedanea) – Wax tree is native of Asia, although it has been planted elsewhere, most notably in Australia and New Zealand. It is a large shrub or tree, up to 8 m tall, somewhat similar to a sumac tree. Because of its beautiful autumn foliage, it has been planted outside of Asia as an ornamental plant, often by gardeners who were apparently unaware of the dangers of allergic reactions. It is now officially classified as a noxious weed in Australia and New Zealand. The fatty-acid methyl ester of the kernel oil meets all of the major biodiesel requirements in the USA (ASTM D 6751-02, ASTM PS 121-99), Germany (DIN V 51606) and European Union (EN 14214). Toxicodendron sylvestre (synonym Rhus sylvestris) – native to China, Japan, Korea and Taiwan. Toxicodendron trichocarpum – southern China, Korea, Japan, and Kuril Islands Toxicodendron vernicifluum (synonym Rhus verniciflua) – Lacquer tree or varnish tree grows in Asia, especially China and Japan. Growing up to 20 m tall, its sap produces an extremely durable lacquer. The leaves have 7–19 leaflets (most often 11–13). The sap contains the allergenic oil, urushiol. Urushiol gets its name from this species which in Japanese is called Urushi. Other names for this species include Japanese lacquer tree, Japanese Varnish Tree, and Japanese Sumac (Note: the term "varnish tree" is also occasionally applied to the Candlenut, Aleurites moluccana, a southeast Asian tree unrelated to Toxicodendron). Toxicodendron vernix (synonym Rhus vernix) – Poison sumac is a tall shrub or a small tree, from 2–7 m tall. It is found in swampy, open areas and reproduces by seeds. The leaves have between 7–13 untoothed leaflets, in a feather-compound arrangement. In terms of its potential to cause urushiol-induced contact dermatitis, poison sumac is far more virulent than other Toxicodendron species, even more virulent than poison ivy and poison oak. According to some botanists, T. vernix is the most toxic plant species in the United States (Frankel, 1991). Toxicodendron wallichii – Himalayas, southern Tibet, southern China, Vietnam, and northern Thailand Toxicodendron yunnanense – Yunnan Formerly placed here Searsia parviflora (as Toxicodendron parviflorum ) – Small-flowered poison sumac grows in the Himalayas between Kumaun, India and Bhutan Toxicity Uses In East Asia, in particular in Japan, traditional candle fuel was produced from Toxicodendron vernicifluum and Toxicodendron succedaneum, among other sumac plants in the genus Toxicodendron, rather than beeswax or animal fats. The sumac wax was a byproduct of traditional Japanese lacquer manufacture. The conical rousoku candles produced from sumac wax burn with smokeless flame and were favored in many respects over candles made from lard or beeswax during the Tokugawa shogunate. Japan wax is not a true wax but a solid fat that contains 10-15% palmitin, stearin, and olein with about 1% japanic acid (1,21-heneicosanedioic acid). It is still used in many tropical and subtropical countries in the production of wax match sticks.	Biology and health sciences	Sapindales	Plants
69232	https://en.wikipedia.org/wiki/Sumac	Sumac	Sumac or sumach ( , )—not to be confused with poison sumac—is any of the roughly 35 species of flowering plants in the genus Rhus (and related genera) of the cashew and mango tree family, Anacardiaceae. However, it is Rhus coriaria that is most commonly used for culinary purposes. Sumac is prized as a spice—especially in Arab, Iranian, Lebanese, and other Middle Eastern cuisines—and used as a dye and holistic remedy. The plants grow in subtropical and temperate regions, on nearly every continent except Antarctica and South America. Description Sumacs are dioecious shrubs and small trees in the family Anacardiaceae that can reach a height of . The leaves are usually pinnately compound, though some species have trifoliate or simple leaves. The flowers are in dense panicles or spikes long, each flower very small, greenish, creamy white or red, with five petals. The fruits are reddish, thin-fleshed drupes covered in varying levels of hairs at maturity and form dense clusters at branch tips, sometimes called sumac bobs. Sumacs propagate both by seed (spread by birds and other animals through their droppings), and by new shoots from rhizomes, forming large clonal colonies. Taxonomy The taxonomy of Rhus has a long history, with de Candolle proposing a subgeneric classification with 5 sections in 1825. At its largest circumscription, Rhus, with over 250 species, has been the largest genus in the family Anacardiaceae. Other authors used subgenera and placed some species in separate genera, hence the use of Rhus sensu lato and Rhus sensu stricto (s.s.). One classification uses two subgenera, Rhus (about 10 spp.) and Lobadium (about 25 spp.), while at the same time Cotinus, Duckera, Malosma, Metopium, Searsia and Toxicodendron segregated to create Rhus s.s.. Other genera that have been segregated include Actinocheita and Baronia. As defined, Rhus s.s. appears monophyletic by molecular phylogeny research. However, the subgenera do not appear to be monophyletic. The larger subgenus, Lobadium, has been divided further into sections, Lobadium, Terebinthifolia, and Styphonia (two subsections). Accepted species by continent As of November 2024, Plants of the World Online accepts 54 species. Asia, North Africa and southern Europe Rhus amherstensis Rhus chinensis Mill. – Chinese sumac Rhus coriaria – Sicilian sumac, Tanner's sumac Rhus dhuna Rhus potaninii – Potanin's lacquer tree or Chinese varnish tree Rhus punjabensis Rhus taishanensis Rhus teniana Rhus wilsonii Australia, Pacific Rhus caudata Rhus lamprocarpa Rhus lenticellosa Rhus linguata Rhus sandwicensis A.Gray – neneleau or Hawaiian sumac (Hawaii) Rhus taitensis Guill. (Northeast Australia, Malesia, Micronesia, French Polynesia) North America Rhus allophyloides Rhus andrieuxii Rhus aromatica – fragrant sumac Rhus arsenei Rhus × ashei (R. glabra × R. michauxii) Rhus bahamensis Rhus barclayi Rhus chondroloma Rhus choriophylla Rhus ciliolata Rhus copallinum – winged or shining sumac Rhus duckerae Rhus galeottii Rhus glabra – smooth sumac Rhus integrifolia – lemonade sumac Rhus jaliscana Rhus kearneyi – Kearney sumac Rhus lanceolata – prairie sumac Rhus lentii Rhus michauxii – Michaux's sumac Rhus microphylla – desert sumac, littleleaf sumac Rhus muelleri Rhus nelsonii Rhus oaxacana Rhus ovata – sugar sumac Rhus pachyrrhachis Rhus palmeri Rhus × pulvinata (R. glabra × R. typhina) Rhus rubifolia Rhus schiedeana Rhus schmidelioides Rhus standleyi Rhus tamaulipana Rhus tepetate Rhus terebinthifolia Rhus trilobata Nutt. – skunkbush sumac Rhus typhina – staghorn sumac Rhus vestita Rhus virens Lindh. ex A.Gray– evergreen sumac †Rhus boothillensis Flynn, DeVore, & Pigg-Ypresian, Washington †Rhus garwellii Flynn, DeVore, & Pigg-Ypresian, Washington †Rhus malloryi Wolfe & Wehr – Ypresian, Washington †Rhus republicensis Flynn, DeVore, & Pigg-Ypresian, Washington †Rhus rooseae Manchester – Middle Eocene, Oregon Formerly placed here Searsia mysorensis (as Rhus mysorensis ) – Mysore sumac Etymology The word sumac traces its etymology from Old French sumac (13th century), from Mediaeval Latin sumach, from Arabic (), from Syriac (ܣܘܡܩܐ)- meaning "red". The generic name Rhus derives from Ancient Greek ῥοῦς (rhous), meaning "sumac", of unknown etymology; the suggestion that it is connected with the verb ῥέω (rheō), "to flow", is now rejected by scholars. Cultivation and uses Species including the fragrant sumac (R. aromatica), the littleleaf sumac (R. microphylla), the smooth sumac (R. glabra), and the staghorn sumac (R. typhina) are grown for ornament, either as the wild types or as cultivars. In food The dried fruits of some species are ground to produce a tangy, crimson spice popular in many countries. Fruits are also used to make a traditional "pink lemonade" beverage by steeping them in water, straining to remove the hairs that may irritate the mouth or throat, sometimes adding sweeteners such as honey or sugar. Sumac's tart flavor comes from high amounts of malic acid. The fruits (drupes) of Rhus coriaria are ground into a reddish-purple powder used as a spice in Middle Eastern cuisine to add a tart, lemony taste to salads or meat. In Arab cuisine, it is used as a garnish on meze dishes such as hummus and tashi, it is also commonly added to falafel. Syria uses the spice also, it is one of the main ingredients of Kubah Sumakieh in Aleppo of Syria, it is added to salads in the Levant, as well as being one of the main ingredients in the Palestinian dish musakhan. In Afghan, Armenian, Iraqi, Iranian and Mizrahi cuisines, sumac is added to rice or kebab. In Armenian, Azerbaijani, Central Asian, Syrian, Iraqi, Jordanian, Palestinian, Lebanese, Turkish and Kurdish cuisines, it is added to salads, kebab and lahmajoun. Rhus coriaria is used in the spice mixture za'atar. During medieval times, primarily from the thirteenth to fifteenth centuries, sumac appeared in cookbooks frequently used by the affluent in Western Europe. One dish in particular called sumāqiyya, a stew made from sumac, was frequently rendered as "somacchia" by Europeans. In North America, the smooth sumac (R. glabra), three-leaf sumac (R. trilobata), and staghorn sumac (R. typhina) are sometimes used to make a beverage termed "sumac-ade", "Indian lemonade", or "rhus juice". This drink is made by soaking the drupes in cool water, rubbing them to extract the essence, straining the liquid through a cotton cloth, and sweetening it. Native Americans also use the leaves and drupes of these sumacs combined with tobacco in traditional smoking mixtures. Dye and tanning agent The leaves and bark of most sumac species contain high levels of tannins and have been used in the manufacturing of leather by many cultures around the world. The Hebrew name og ha-bursaka'im means "tanner's sumac", as does the Latin name of R. coriaria. The leaves of certain sumacs yield tannin (mostly pyrogallol-type), a substance used in vegetable tanning. Notable sources include the leaves of R. coriaria, Chinese gall on R. chinensis, and wood and roots of R. pentaphylla. Leather tanned with sumac is flexible, light in weight, and light in color. One type of leather made with sumac tannins is morocco leather. The dyeing property of sumac needed to be considered when it was shipped as a fine floury substance in sacks as a light cargo accompanying heavy cargoes such as marble. Sumac was especially dangerous to marble: "When sumac dust settles on white marble, the result is not immediately apparent, but if it once becomes wet, or even damp, it becomes a powerful purple dye, which penetrates the marble to an extraordinary depth." Ibn Badis describes a formula for making red ink out of leeched sumac mixed with gum. was used only for the outerwear of the Emperor of Japan, thus being one of the forbidden сolors. Traditional medicinal use Sumac was used as a treatment for several different ailments in medieval medicine, primarily in Middle Eastern and South Asian countries (where sumac was more readily available than in Europe). An 11th-century shipwreck off the coast of Rhodes, excavated by archeologists in the 1970s, contained commercial quantities of sumac drupes. These could have been intended for use as medicine, as a culinary spice, or as a dye. A clinical study showed that dietary sumac decreases the blood pressure in patients with hypertension and can be used as adjunctive treatment. Other uses Some beekeepers use dried sumac bobs as a source of fuel for their smokers. Sumac stems also have a soft pith in the center that is easily removed to make them useful in traditional Native American pipemaking. They were commonly used as pipe stems in the northern United States. Dried sumac wood fluoresces under long-wave ultraviolet radiation. Toxicity and control Some species formerly recognized in Rhus, such as poison ivy (Toxicodendron radicans, syn. Rhus toxicodendron), poison oak (Toxicodendron diversilobum, syn. Rhus diversiloba), and poison sumac (Toxicodendron vernix, syn. Rhus vernix), produce the allergen urushiol and can cause severe delayed hypersensitivity reactions. Poison sumac may be identified by its white drupes, which are quite different from the red drupes of true Rhus species. Mowing of sumac is not a good control measure, since the wood is springy, resulting in jagged, sharp-pointed stumps when mown. The plant will quickly recover with new growth after mowing. Goats have long been considered an efficient and quick removal method, as they eat the bark, which helps prevent new shoots. Sumac propagates by rhizome. Small shoots will be found growing near a more mature sumac tree via a shallow running root quite some distance from the primary tree. Thus, root pruning is a means of control without eliminating the plants altogether. Explanatory notes	Biology and health sciences	Herbs and spices	Plants
69279	https://en.wikipedia.org/wiki/Angioplasty	Angioplasty	Angioplasty, also known as balloon angioplasty and percutaneous transluminal angioplasty (PTA), is a minimally invasive endovascular procedure used to widen narrowed or obstructed arteries or veins, typically to treat arterial atherosclerosis. A deflated balloon attached to a catheter (a balloon catheter) is passed over a guide-wire into the narrowed vessel and then inflated to a fixed size. The balloon forces expansion of the blood vessel and the surrounding muscular wall, allowing an improved blood flow. A stent may be inserted at the time of ballooning to ensure the vessel remains open, and the balloon is then deflated and withdrawn. Angioplasty has come to include all manner of vascular interventions that are typically performed percutaneously. Uses and indications Coronary angioplasty A coronary angioplasty is a therapeutic procedure to treat the stenotic (narrowed) coronary arteries of the heart found in coronary heart disease. These stenotic segments of the coronary arteries arise due to the buildup of cholesterol-laden plaques that form in a condition known as atherosclerosis. A percutaneous coronary intervention (PCI), or coronary angioplasty with stenting, is a non-surgical procedure used to improve the blood flow to the heart. Coronary angioplasty is indicated for coronary artery diseases such as unstable angina, NSTEMI, STEMI and spontaneous coronary artery perforation. PCI for stable coronary disease has been shown to significantly relieve symptoms such as angina, or chest pain, thereby improving functional limitations and quality of life. Peripheral angioplasty Peripheral angioplasty refers to the use of a balloon to open a blood vessel outside the coronary arteries. It is most commonly done to treat atherosclerotic narrowings of the abdomen, leg and renal arteries caused by peripheral artery disease. Often, peripheral angioplasty is used in conjunction with guide wire, peripheral stenting and an atherectomy. Chronic limb-threatening ischemia Angioplasty can be used to treat advanced peripheral artery disease to relieve the claudication, or leg pain, that is classically associated with the condition. The bypass versus angioplasty in severe ischemia of the leg (BASIL) trial investigated infrainguinal bypass surgery first compared to angioplasty first in select patients with severe lower limb ischemia who were candidates for either procedure. The BASIL trial found that angioplasty was associated with less short term morbidity compared with bypass surgery, however long term outcomes favor bypass surgery. Based on the BASIL trial, the ACCF/AHA guidelines recommend balloon angioplasty only for patients with a life expectancy of 2 years or less or those who do not have an autogenous vein available. For patients with a life expectancy greater than 2 of years life, or who have an autogenous vein, a bypass surgery could be performed first. Renal artery angioplasty Renal artery stenosis is associated with hypertension and loss of renal function. Atherosclerotic obstruction of the renal artery can be treated with angioplasty with or without stenting of the renal artery. There is a weak recommendation for renal artery angioplasty in patients with renal artery stenosis and flash edema or congestive heart failure. Carotid angioplasty Carotid artery stenosis can be treated with angioplasty and carotid stenting for patients at high risk for undergoing carotid endarterectomy (CEA). Although carotid endarterectomy is typically preferred over carotid artery stenting, stenting is indicated in select patients with radiation-induced stenosis or a carotid lesion not suitable for surgery. Venous angioplasty Angioplasty is used to treat venous stenosis affecting dialysis access, with drug-coated balloon angioplasty proving to have better 6 month and 12 month patency than conventional balloon angioplasty. Angioplasty is occasionally used to treat residual subclavian vein stenosis following decompression surgery for thoracic outlet syndrome. There is a weak recommendation for deep venous stenting to treat obstructive chronic venous disease. Contraindications Angioplasty requires an access vessel, typically the femoral or radial artery or femoral vein, to permit access to the vascular system for the wires and catheters used. If no access vessel of sufficient size and quality is available, angioplasty is contraindicated. A small vessel diameter, the presence of posterior calcification, occlusion, hematoma, or an earlier placement of a bypass origin, may make access to the vascular system too difficult. Percutaneous transluminal coronary angioplasty (PTCA) is contraindicated in patients with left main coronary artery disease, due to the risk of spasm of the left main coronary artery during the procedure. Also, PTCA is not recommended if there is less than 70% stenosis of the coronary arteries, as the stenosis is not deemed to be hemodynamically significant below this level. Technique Access to the vascular system is typically gained percutaneously (through the skin, without a large surgical incision). An introducer sheath is inserted into the blood vessel via the Seldinger technique. Fluoroscopic guidance uses magnetic resonance or X-ray fluoroscopy and radiopaque contrast dye to guide angled wires and catheters to the region of the body to be treated in real time. Tapered guidewire is chosen for small occlusion, followed by intermediate type guidewires for tortuous arteries and difficulty passing through extremely narrow channels, and stiff wires for hard, dense, and blunt occlusions. To treat a narrowing in a blood vessel, a wire is passed through the stenosis in the vessel and a balloon on a catheter is passed over the wire and into the desired position. The positioning is verified by fluoroscopy and the balloon is inflated using water mixed with contrast dye to 75 to 500 times normal blood pressure (6 to 20 atmospheres), with most coronary angioplasties requiring less than 10 atmospheres. A stent may or may not also be placed. At the conclusion of the procedure, the balloons, wires and catheters are removed and the vessel puncture site is treated either with direct pressure or a vascular closure device. Transradial artery access (TRA) and transfemoral artery access (TFA) are two techniques for percutaneous coronary intervention. TRA is the technique of choice for management of acute coronary syndrome (ACS) as it has significantly lower incidence of bleeding and vascular complications compared with the TFA approach. TRA also has a mortality benefit for high risk ACS patients and high risk bleeding patients. TRA was also found to yield improved quality of life, as well as decreased healthcare costs and resources. Risks and complications Relative to surgery, angioplasty is a lower-risk option for the treatment of the conditions for which it is used, but there are unique and potentially dangerous risks and complications associated with angioplasty: Embolization, or the launching of debris into the bloodstream Bleeding from over-inflation of a balloon catheter or the use of an inappropriately large or stiff balloon, or the presence of a calcified target vessel. Hematoma or pseudoaneurysm formation at the access site Radiation-induced injuries (burns) from the X-rays used Contrast-induced renal injury Cerebral Hyperperfusion Syndrome leading to stroke is a serious complication of carotid artery angioplasty with stenting. Angioplasty may also provide a less durable treatment for atherosclerosis and be more prone to restenosis relative to vascular bypass or coronary artery bypass grafting. Drug-eluting balloon angioplasty has significantly less restenosis, late lumen loss and target lesion revascularization at both short term and midterm follow-up compared to uncoated balloon angioplasty for femoropopliteal arterial occlusive disease. Although angioplasty of the femoropopliteal artery with paclitaxel-coated stents and balloons significantly reduces rates of vessel restenosis and target lesion revascularization, it was also found to have increased risk of death. Recovery After angioplasty, most patients are monitored overnight in the hospital, but if there are no complications, patients are sent home the following day. The catheter site is checked for bleeding and swelling and the heart rate and blood pressure are monitored to detect late rupture and hemorrhage. Post-procedure protocol also involves monitoring urinary output, cardiac symptoms, pain and other signs of systemic problems. Usually, patients receive medication that will relax them to protect the arteries against spasms. Patients are typically able to walk within two to six hours following the procedure and return to their normal routine by the following week. Angioplasty recovery consists of avoiding physical activity for several days after the procedure. Patients are advised to avoid heavy lifting and strenuous activities for a week. Patients will need to avoid physical stress or prolonged sport activities for a maximum of two weeks after a delicate balloon angioplasty. After the initial two week recovery phase, most angioplasty patients can begin to safely return to low-level exercise. A graduated exercise program is recommended whereby patients initially perform several short bouts of exercise each day, progressively increasing to one or two longer bouts of exercise. As a precaution, all structured exercise should be cleared by a cardiologist before commencing. Exercise-based rehabilitation following percutaneous coronary intervention has shown improvement in recurrent angina, total exercise time, ST-segment decline, and maximum exercise tolerance. Patients who experience swelling, bleeding or pain at the insertion site, develop fever, feel faint or weak, notice a change in temperature or color in the arm or leg that was used or have shortness of breath or chest pain should immediately seek medical advice. Patients with stents are usually prescribed dual antiplatelet therapy (DAPT) which consists of a P2Y12 inhibitor, such as clopidogrel, which is taken at the same time as acetylsalicylic acid (aspirin). Dual antiplatelet therapy (DAPT) is recommended for 1 month following bare metal stent placement, for 3 months following a second generation drug-eluting stent placement, and for 6–12 months following a first generation drug-eluting stent placement. DAPT's antiplatelet properties are intended to prevent blood clots, however they also increase the risk of bleeding, so it is important to consider each patient's preferences, cardiac conditions, and bleeding risk when determining the duration of DAPT treatment. Another important consideration is that concomitant use of Clopidogrel and Proton Pump Inhibitors following coronary angiography is associated with significantly higher adverse cardiovascular complications such as major adverse cardiovascular events (MACE), stent thrombosis and myocardial infarction. History Angioplasty was first described by the US interventional radiologist Charles Dotter in 1964. Dotter pioneered modern medicine with the invention of angioplasty and the catheter-delivered stent, which were first used to treat peripheral arterial disease. On January 16, 1964, Dotter percutaneously dilated a tight, localized stenosis of the subsartorial artery in an 82-year-old woman with painful leg ischemia and gangrene who refused leg amputation. After successful dilation of the stenosis with a guide wire and coaxial Teflon catheters, the circulation returned to her leg. The dilated artery stayed open until her death from pneumonia two and a half years later. Charles Dotter is commonly known as the "Father of Interventional Radiology" and was nominated for the Nobel Prize in medicine in 1978. The first percutaneous coronary angioplasty on an awake patient was performed in Zurich by the German cardiologist Andreas Gruentzig on September 16, 1977. The first percutaneous coronary angioplasties in the United States were performed on the same day (March 1, 1978) by Simon H. Stertzer at Lenox Hill Hospital in New York and Richard K. Myler at St. Mary's Hospital in San Francisco. During the previous year, also at St. Mary's Hospital in San Francisco, Myler and Gruentzig had performed dilatations in the setting of bypass surgery to test the catheter concept before Gruentzig performed the first PTCA in his catheterization lab in Zurich. The initial form of angioplasty was 'plain old balloon angioplasty' (POBA) without stenting, until the invention of bare metal stenting in the mid-1980s to prevent the abrupt closure sometimes seen with POBA. Bare metal stents were found to cause in-stent restenosis as a result of neointimal hyperplasia and stent thrombosis, which led to the invention of drug-eluting stents with anti-proliferative drugs to combat in-stent restenosis. The first coronary angioplasty with a drug delivery stent system was performed by Stertzer and Luis de la Fuente, at the Instituto Argentino de Diagnóstico y Tratamiento (English: Argentina Institute of Diagnosis and Treatment) in Buenos Aires, in 1999. Ingemar Henry Lundquist invented the over-the-wire balloon catheter that is now used in the majority of angioplasty procedures in the world. A subset of angioplasty, known as excimer laser coronary angioplasty (ELCA), uses excimer lasers to remove small amounts of tissue, including undilatable and uncrossable lesions, in the artery in order to allow the balloon to more effectively compress plaque into the artery walls. Such work was first developed in 1984 following earlier work in 1980–1983, when Rangaswamy Srinivasan, Samuel Blum and James J. Wynne at IBM's T. J. Watson Research Center observed the effect of the ultraviolet excimer laser on biological materials. Intrigued, they investigated further, finding that the laser made clean, precise cuts that would be ideal for delicate surgeries. This resulted in a fundamental patent and Srinivasan, Blum and Wynne were elected to the National Inventors Hall of Fame in 2002. In 2012, the team members were honored with National Medal of Technology and Innovation by the President Barack Obama for their work related to the excimer laser. Robert Ginsburg deployed the first used of ELCA in 1984 on a patient with severe stenosis of the deep femoral artery and a threatened limb.	Biology and health sciences	Surgery	Health
69398	https://en.wikipedia.org/wiki/Delta%20Works	Delta Works	The Delta Works () is a series of construction projects in the southwest of the Netherlands to protect a large area of land around the Rhine–Meuse–Scheldt delta from the sea. Constructed between 1954 and 1997, the works consist of dams, sluices, locks, dykes, levees, and storm surge barriers located in the provinces of South Holland and Zeeland. The aim of the dams, sluices, and storm surge barriers was to shorten the Dutch coastline, thus reducing the number of dikes that had to be raised. Along with the Zuiderzee Works, the Delta Works have been declared one of the Seven Wonders of the Modern World by the American Society of Civil Engineers. History Due to indecision and the Second World War, little action was taken. In 1950 two small estuary mouths, the Brielse Gat near Brielle and the Botlek near Vlaardingen were dammed. After the North Sea flood of 1953, a Delta Works Commission was installed to research the causes and develop measures to prevent such disasters in future. They revised some of the old plans and came up with the "Deltaplan". Delta law and conceptual framework An important part of this project was fundamental research to come up with long-term solutions, protecting the Netherlands against future floods. Instead of analysing past floods and building protection sufficient to deal with those, the Delta Works commission pioneered a conceptual framework to use as norm for investment in flood defences. The framework is called the 'Delta norm'; it includes the following principles: Major areas to be protected from flooding are identified. These are called "dike ring areas" because they are protected by a ring of primary sea defences. The cost of flooding is assessed using a statistical model involving damage to property, lost production, and a given amount per human life lost. For the purpose of this model, a human life is valued at €2.2 million (2008 data). The chances of a significant flood within the given area are calculated. This is done using data from a purpose-built flood simulation lab, as well as empirical statistical data regarding water wave properties and distribution. Storm behaviour and spring tide distribution are also taken into account. The most important "dike ring area" is the South Holland coast region. It is home to four million people, most of whom live below normal sea level. The loss of human life in a catastrophic flood here can be very large because there is typically little warning time with North Sea storms. Comprehensive evacuation is not a realistic option for the Holland coastal region. The commission initially set the acceptable risk for complete failure of every "dike ring" in the country at 1 in 125,000 years. But, it found that the cost of building to this level of protection could not be supported. It set "acceptable" risks by region as follows: North and South Holland (excluding Wieringermeer): 1 per 10,000 years Other areas at risk from sea flooding: 1 per 4,000 years Transition areas between high land and low land: 1 per 2,000 years River flooding causes less damage than salt water flooding, which causes long-term damage to agricultural lands. Areas at risk from river flooding were assigned a higher acceptable risk. River flooding also has a longer warning time, producing a lower estimated death toll per event. South Holland at risk from river flooding: 1 per 1,250 years Other areas at risk from river flooding: 1 per 250 years. These acceptable risks were enshrined in the Delta Law (Dutch: Deltawet). This required the government to keep risks of catastrophic flooding within these limits and to upgrade defences should new insights into risks require this. The limits have also been incorporated into the new Water Law (Waterwet), effective from 22 December 2009. The Delta Project (of which the Delta Works are a part) has been designed with these guidelines in mind. All other primary defences have been upgraded to meet the norm. New data elevating the risk assessment on expected sea level rise due to global warming has identified ten 'weak points.' These have been upgraded to meet future demands. The latest upgrades are made under the High Water Protection Program. Alterations to the plan during the execution of the Works During the execution of the works, changes were made in response to public pressure. In the Nieuwe Waterweg, the heightening and the associated widening of the dikes proved very difficult because of public opposition to the planned destruction of important historic buildings to achieve this. The plan was changed to the construction of a storm surge barrier (the Maeslantkering) and dikes were only partly built up. The storm-surge barrier The Delta Plan originally intended to create a large freshwater lake, the (Zeeland Lake). This would have caused major environmental destruction in Oosterschelde, with the total loss of the saltwater ecosystem and, consequently, the harvesting of oysters. Environmentalists and fishermen combined their efforts to prevent the closure; they persuaded parliament to amend the original plan. Instead of completely damming the estuary, the government agreed to build a storm surge barrier. This essentially is a long collection of very large valves that can be closed against storm surges. The storm surge barrier closes only when the sea-level is expected to rise 3 metres above mean sea level. Under normal conditions, the estuary's mouth is open, and salt water flows in and out with the tide. As a result of the change, the weak dikes along the Oosterschelde needed to be strengthened. Over 200 km of the dike needed new revetments. The connections between the Eastern Scheldt and the neighboring Haringvliet had to be dammed to limit the effect of the salt water. Extra dams and locks were needed at the east part of the Oosterschelde to create a shipping route between the ports of Rotterdam and Antwerp. Since operating the barrier has an effect on the environment, fisheries and the water management system, decisions made on opening or closing the gate are carefully considered. Also the safety of the surrounding dykes are affected by barrier operations. Environmental policy implementations In an attempt to restore and preserve the natural system surrounded by the dykes and storm-surge barrier, the concept 'building with nature' was introduced in revised Delta Program updates after 2008. The new integrated water management plan not only takes into account protection against flooding, but also covers water quality, leisure industry, economic activities, shipping, environment and nature. Whenever possible, existing engineering constructions would be replaced by more 'nature friendly' options in an attempt to restore natural estuary and tides, while still protecting against flooding. In addition, building components of the reinforcements are designed in a way that they support formation of entire ecosystems. As part of the revision, the Room for the River projects, enabled nature to occupy space by lowering or widening the river bed. In order to establish this, agricultural flood plains are turned into natural parks, excavated farmland is used for wild vegetation and newly excavated lakes and bypasses create habitats for fish and birds. Along the coast, natural sand is added each year to allow sand to blow freely through the dunes instead of having the dunes held in place by planted vegetation or revetments. Although the new plan brought along additional cost, it was received favourably. The re-considerations of the Delta Project indicated the growing importance of integrate environmental impact assessments in policy-making. Environmental effects The Delta Project of which the Delta Works are part of was originally designed in a period of time when environmental awareness and ecological effects of engineering projects were barely taken into consideration. Although the level of awareness for the environment grew throughout the years, the Delta Project has caused numerous irreversible effects on the environment in the past. Blocking the estuary mouths did reduce the length of dykes that otherwise would have to be built to protect against floods, but it also led to major changes in the water systems. For example, the tides disappeared, which resulted in a less smooth transition from sea water into fresh water. Flora and fauna suffered from this noticeable change. In addition, rivers experienced an increase in sedimentation by polluted sludge, since there was no longer an open passage to the sea. Project costs The projects of the Delta Plan are financed with the Delta Fund. In 1958, when the Delta law was accepted under the Delta Works Commission, the total costs were estimated at 3.3 billion guilder. This was at that time equal to 20% of national GDP. This amount was spread out over the 25 years that it would take to complete the massive engineering project. The Delta works were mostly financed by the national budget, with a contribution of the Marshall Plan of 400 million guilder. In addition, the Dutch natural gas discovery contributed massively to the finance of the project. At completion in 1997, costs were set on 8.2 billion guilder. Nevertheless, in 2012 the total costs were already set on around $13 billion. Current status The original plan was completed by the Europoortkering which required the construction of the Maeslantkering in the Nieuwe Waterweg between Maassluis and Hook of Holland and the Hartelkering in the Hartel Canal near Spijkenisse. The works were declared finished after almost forty years in 1997. Due to climate change and relative sea-level rise, the dikes will eventually have to be made higher and wider. This is a long term uphill battle against the sea. The needed level of flood protection and the resulting costs are a recurring subject of debate, and involve a complicated decision-making process. In 1995 it was agreed in the Delta Plan Large Rivers and Room for the River projects that about 500 kilometres of insufficient dyke revetments were reinforced and replaced along the Oosterschelde and Westerschelde between 1995 and 2015. After 2015, under the High Water Protection Program, additional upgrades are made. In September 2008, the Delta Commission presided by politician Cees Veerman advised in a report that the Netherlands would need a massive new building program to strengthen the country's water defenses against the anticipated effects of global warming for the next 190 years. The plans included drawing up worst-case scenarios for evacuations and included more than €100 billion, or $144 billion, in new spending through the year 2100 for measures, such as broadening coastal dunes and strengthening sea and river dikes. The commission said the country must plan for a rise in the North Sea of 1.3 meters by 2100 and 4 meters by 2200. Projects The works that are part of the Delta Works are listed in chronological order with their year of completion:	Technology	Dams	null
69427	https://en.wikipedia.org/wiki/Botanical%20garden	Botanical garden	A botanical garden or botanic garden is a garden with a documented collection of living plants for the purpose of scientific research, conservation, display, and education. It is their mandate as a botanical garden that plants are labelled with their botanical names. It may contain specialist plant collections such as cacti and other succulent plants, herb gardens, plants from particular parts of the world, and so on; there may be glasshouses or shadehouses, again with special collections such as tropical plants, alpine plants, or other exotic plants that are not native to that region. Most are at least partly open to the public, and may offer guided tours, public programming such as workshops, courses, educational displays, art exhibitions, book rooms, open-air theatrical and musical performances, and other entertainment. Botanical gardens are often run by universities or other scientific research organizations, and often have associated herbaria and research programmes in plant taxonomy or some other aspect of botanical science. In principle, their role is to maintain documented collections of living plants for the purposes of scientific research, conservation, display, and education, although this will depend on the resources available and the special interests pursued at each particular garden. The staff will normally include botanists as well as gardeners. Many botanical gardens offer diploma/certificate programs in horticulture, botany and taxonomy. There are many internship opportunities offered to aspiring horticulturists. As well as opportunities for students/researchers to use the collection for their studies. History The origin of modern botanical gardens is generally traced to the appointment of botany professors to the medical faculties of universities in 16th-century Renaissance Italy, which also entailed curating a medicinal garden. However, the objectives, content, and audience of today's botanic gardens more closely resembles that of the grandiose gardens of antiquity and the educational garden of Theophrastus in the Lyceum of ancient Athens. The early concern with medicinal plants changed in the 17th century to an interest in the new plant imports from explorations outside Europe as botany gradually established its independence from medicine. In the 18th century, systems of nomenclature and classification were devised by botanists working in the herbaria and universities associated with the gardens, these systems often being displayed in the gardens as educational "order beds". With the rapid expansion of European colonies around the globe in the late 18th century, botanic gardens were established in the tropics, and economic botany became a focus with the hub at the Royal Botanic Gardens, Kew, near London. Over the years, botanical gardens, as cultural and scientific organisations, have responded to the interests of botany and horticulture. Nowadays, most botanical gardens display a mix of the themes mentioned and more; having a strong connection with the general public, there is the opportunity to provide visitors with information relating to the environmental issues being faced at the start of the 21st century, especially those relating to plant conservation and sustainability. Definitions The "New Royal Horticultural Society Dictionary of Gardening" (1999) points out that among the various kinds of organizations known as botanical gardens, there are many that are in modern times public gardens with little scientific activity, and it cited a tighter definition published by the World Wildlife Fund and IUCN when launching the "Botanic Gardens Conservation Strategy" in 1989: "A botanic garden is a garden containing scientifically ordered and maintained collections of plants, usually documented and labelled, and open to the public for the purposes of recreation, education and research." The term tends to be used somewhat differently in different parts of the world. For example a large woodland garden with a good collection of rhododendron and other flowering tree and shrub species is very likely to present itself as a "botanical garden" if it is located in the US, but very unlikely to do so if in the UK (unless it also contains other relevant features). Very few of the sites used for the UK's dispersed National Plant Collection, usually holding large collections of a particular taxonomic group, would call themselves "botanic gardens". This has been further reduced by Botanic Gardens Conservation International to the following definition which "encompasses the spirit of a true botanic garden": "A botanic garden is an institution holding documented collections of living plants for the purposes of scientific research, conservation, display and education." The following definition was produced by staff of the Liberty Hyde Bailey Hortorium of Cornell University in 1976. It covers in some detail the many functions and activities generally associated with botanical gardens: A botanical garden is a controlled and staffed institution for the maintenance of a living collection of plants under scientific management for purposes of education and research, together with such libraries, herbaria, laboratories, and museums as are essential to its particular undertakings. Each botanical garden naturally develops its own special fields of interests depending on its personnel, location, extent, available funds, and the terms of its charter. It may include greenhouses, test grounds, an herbarium, an arboretum, and other departments. It maintains a scientific as well as a plant-growing staff, and publication is one of its major modes of expression. This broad outline is then expanded: The botanic garden may be an independent institution, a governmental operation, or affiliated to a college or university. If a department of an educational institution, it may be related to a teaching program. In any case, it exists for scientific ends and is not to be restricted or diverted by other demands. It is not merely a landscaped or ornamental garden, although it may be artistic, nor is it an experiment station or yet a park with labels on the plants. The essential element is the intention of the enterprise, which is the acquisition and dissemination of botanical knowledge. A contemporary botanic garden is a strictly protected green area, where a managing organization creates landscaped gardens and holds documented collections of living plants and/or preserved plant accessions containing functional units of heredity of actual or potential value for purposes such as scientific research, education, public display, conservation, sustainable use, tourism and recreational activities, production of marketable plant-based products and services for improvement of human well-being. The botanical gardens network Worldwide, there are now about 1800 botanical gardens and arboreta in about 150 countries (mostly in temperate regions) of which about 550 are in Europe (150 of which are in Russia), 200 in North America, and an increasing number in East Asia. These gardens attract about 300 million visitors a year. Historically, botanical gardens exchanged plants through the publication of seed lists (these were called in the 18th century). This was a means of transferring both plants and information between botanical gardens. This system continues today, although the possibility of genetic piracy and the transmission of invasive species has received greater attention in recent times. The International Association of Botanic Gardens was formed in 1954 as a worldwide organisation affiliated to the International Union of Biological Sciences. More recently, coordination has also been provided by Botanic Gardens Conservation International (BGCI), which has the mission "To mobilise botanic gardens and engage partners in securing plant diversity for the well-being of people and the planet". BGCI has over 700 membersmostly botanic gardensin 118 countries, and strongly supports the Global Strategy for Plant Conservation by producing a range resources and publications, and by organizing international conferences and conservation programs. Communication also happens regionally. In the United States, there is the American Public Gardens Association (formerly the American Association of Botanic Gardens and Arboreta), and in Australasia there is the Botanic Gardens of Australia and New Zealand (BGANZ). History and development The history of botanical gardens is closely linked to the history of botany itself. The botanical gardens of the 16th and 17th centuries were medicinal gardens, but the idea of a botanical garden changed to encompass displays of the beautiful, strange, new and sometimes economically important plant trophies being returned from the European colonies and other distant lands. Later, in the 18th century, they became more educational in function, demonstrating the latest plant classification systems devised by botanists working in the associated herbaria as they tried to order these new treasures. Then, in the 19th and 20th centuries, the trend was towards a combination of specialist and eclectic collections demonstrating many aspects of both horticulture and botany. Precursors The idea of "scientific" gardens used specifically for the study of plants dates back to antiquity. Grand gardens of ancient history Near-eastern royal gardens set aside for economic use or display and containing at least some plants gained by special collecting trips or military campaigns abroad, are known from the second millennium BCE in ancient Egypt, Mesopotamia, Crete, Mexico and China. In about 2800 BCE, the Chinese Emperor Shen Nung sent collectors to distant regions searching for plants with economic or medicinal value. It has also been suggested that the Spanish colonization of Mesoamerica influenced the history of the botanical garden as gardens in Tenochtitlan established by king Nezahualcoyotl, also gardens in Chalco (altépetl) and elsewhere, greatly impressed the Spanish invaders, not only with their appearance, but also because the indigenous Aztecs employed many more medicinal plants than did the classical world of Europe. Early medieval gardens in Islamic Spain resembled botanic gardens of the future, an example being the 11th-century Huerta del Rey garden of physician and author Ibn Wafid (999–1075 CE) in Toledo. This was later taken over by garden chronicler Ibn Bassal (fl. 1085 CE) until the Christian conquest in 1085 CE. Ibn Bassal then founded a garden in Seville, most of its plants being collected on a botanical expedition that included Morocco, Persia, Sicily, and Egypt. The medical school of Montpelier was also founded by Spanish Arab physicians, and by 1250 CE, it included a physic garden, but the site was not given botanic garden status until 1593. Physic gardens Botanical gardens, in the modern sense, developed from physic gardens, whose main purpose was to cultivate herbs for medical use as well as research and experimentation. Such gardens have a long history. In Europe, for example, Aristotle (384 BCE – 322 BCE) is said to have had a physic garden in the Lyceum at Athens, which was used for educational purposes and for the study of botany, and this was inherited, or possibly set up, by his pupil Theophrastus, the "Father of Botany". There is some debate among science historians whether this garden was ordered and scientific enough to be considered "botanical", and suggest it more appropriate to attribute the earliest known botanical garden in Europe to the botanist and pharmacologist Antonius Castor, mentioned by Pliny the Elder in the 1st century. Though these ancient gardens shared some of the characteristics of present-day botanical gardens, the forerunners of modern botanical gardens are generally regarded as being the medieval monastic physic gardens that originated after the decline of the Roman Empire at the time of Emperor Charlemagne (742–789 CE). These contained a , a garden used mostly for vegetables, and another section set aside for specially labelled medicinal plants and this was called the or more generally known as a physic garden, and a or orchard. These gardens were probably given impetus when Charlemagne issued a capitulary, the Capitulary de Villis, which listed 73 herbs to be used in the physic gardens of his dominions. Many of these were found in British gardens even though they only occurred naturally in continental Europe, demonstrating earlier plant introduction. Pope Nicholas V set aside part of the Vatican grounds in 1447, for a garden of medicinal plants that were used to promote the teaching of botany, and this was a forerunner to the University gardens at Padua and Pisa established in the 1540s. Certainly the founding of many early botanic gardens was instigated by members of the medical profession. 16th- and 17th-century European gardens In the 17th century, botanical gardens began their contribution to a deeper scientific curiosity about plants. If a botanical garden is defined by its scientific or academic connection, then the first true botanical gardens were established with the revival of learning that occurred in the European Renaissance. These were secular gardens attached to universities and medical schools, used as resources for teaching and research. The superintendents of these gardens were often professors of botany with international reputations, a factor that probably contributed to the creation of botany as an independent discipline rather than a descriptive adjunct to medicine. Origins in the Italian Renaissance The botanical gardens of Southern Europe were associated with university faculties of medicine and were founded in Italy at Orto botanico di Pisa (1544), Orto botanico di Padova (1545), Orto Botanico di Firenze (1545), Orto Botanico dell'Università di Pavia (1558) and Orto Botanico dell'Università di Bologna (1568). Here the physicians (referred to in English as apothecaries) delivered lectures on the Mediterranean "simples" or "officinals" that were being cultivated in the grounds. Student education was no doubt stimulated by the relatively recent advent of printing and the publication of the first herbals. All of these botanical gardens still exist, mostly in their original locations. Northern Europe The tradition of these Italian gardens passed into Spain Botanical Garden of Valencia, 1567) and Northern Europe, where similar gardens were established in the Netherlands (Hortus Botanicus Leiden, 1590; Hortus Botanicus (Amsterdam), 1638), Germany (Alter Botanischer Garten Tübingen, 1535; Leipzig Botanical Garden, 1580; Botanischer Garten Jena, 1586; Botanischer Garten Heidelberg, 1593; Herrenhäuser Gärten, Hanover, 1666; Botanischer Garten der Christian-Albrechts-Universität zu Kiel, 1669; Botanical Garden in Berlin, 1672), Switzerland (Old Botanical Garden, Zürich, 1560; Basel, 1589); England (University of Oxford Botanic Garden, 1621; Chelsea Physic Garden, 1673); Scotland (Royal Botanic Garden Edinburgh, 1670); and in France (Jardin des plantes de Montpellier, 1593; Faculty of Medicine Garden, Paris, 1597; Jardin des Plantes, Paris, 1635), Denmark (University of Copenhagen Botanical Garden, 1600); Sweden (Uppsala University, 1655). Beginnings of botanical science During the 16th and 17th centuries, the first plants were being imported to these major Western European gardens from Eastern Europe and nearby Asia (which provided many bulbs), and these found a place in the new gardens, where they could be conveniently studied by the plant experts of the day. For example, Asian introductions were described by Carolus Clusius (1526–1609), who was director, in turn, of the Botanical Garden of the University of Vienna and Hortus Botanicus Leiden. Many plants were being collected from the Near East, especially bulbous plants from Turkey. Clusius laid the foundations of Dutch tulip breeding and the bulb industry, and he helped create one of the earliest formal botanical gardens of Europe at Leyden where his detailed planting lists have made it possible to recreate this garden near its original site. The of Leyden in 1601 was a perfect square divided into quarters for the four continents, but by 1720, though, it was a rambling system of beds, struggling to contain the novelties rushing in, and it became better known as the . His Exoticorum libri decem (1605) is an important survey of exotic plants and animals that is still consulted today. The inclusion of new plant introductions in botanic gardens meant their scientific role was now widening, as botany gradually asserted its independence from medicine. In the mid to late 17th century, the Paris Jardin des Plantes was a centre of interest with the greatest number of new introductions to attract the public. In England, the Chelsea Physic Garden was founded in 1673 as the "Garden of the Society of Apothecaries". The Chelsea garden had heated greenhouses, and in 1723 appointed Philip Miller (1691–1771) as head gardener. He had a wide influence on both botany and horticulture, as plants poured into it from around the world. The garden's golden age came in the 18th century, when it became the world's most richly stocked botanical garden. Its seed-exchange programme was established in 1682 and still continues today. 18th century With the increase in maritime trade, ever more plants were being brought back to Europe as trophies from distant lands, and these were triumphantly displayed in the private estates of the wealthy, in commercial nurseries, and in the public botanical gardens. Heated conservatories called "orangeries", such as the one at Kew, became a feature of many botanical gardens. Industrial expansion in Europe and North America resulted in new building skills, so plants sensitive to cold were kept over winter in progressively elaborate and expensive heated conservatories and glasshouses. The Cape, Dutch East Indies The 18th century was marked by introductions from the Cape of South Africaincluding ericas, geraniums, pelargoniums, succulents, and proteaceous plantswhile the Dutch trade with the Dutch East Indies resulted in a golden era for the Leiden and Amsterdam botanical gardens and a boom in the construction of conservatories. Royal Botanic Gardens, Kew The Royal Gardens at Kew were founded in 1759, initially as part of the Royal Garden set aside as a physic garden. William Aiton (1741–1793), the first curator, was taught by garden chronicler Philip Miller of the Chelsea Physic Garden whose son Charles became first curator of the original Cambridge Botanic Garden (1762). In 1759, the "Physick Garden" was planted, and by 1767, it was claimed that "the Exotick Garden is by far the richest in Europe". Gardens such as the Royal Botanic Gardens, Kew (1759) and Orotava Acclimatization Garden , Tenerife (1788) and the Real Jardín Botánico de Madrid (1755) were set up to cultivate new species returned from expeditions to the tropics; they also helped found new tropical botanical gardens. From the 1770s, following the example of the French and Spanish, amateur collectors were supplemented by official horticultural and botanical plant hunters. These botanical gardens were boosted by the flora being sent back to Europe from various European colonies around the globe. At this time, British horticulturalists were importing many woody plants from Britain's colonies in North America, and the popularity of horticulture had increased enormously, encouraged by the horticultural and botanical collecting expeditions overseas fostered by the directorship of Sir William Jackson Hooker and his keen interest in economic botany. At the end of the 18th century, Kew, under the directorship of Sir Joseph Banks, enjoyed a golden age of plant hunting, sending out collectors to the South African Cape, Australia, Chile, China, Ceylon, Brazil, and elsewhere, and acting as "the great botanical exchange house of the British Empire". From its earliest days to the present, Kew has in many ways exemplified botanic garden ideals, and is respected worldwide for the published work of its scientists, the education of horticultural students, its public programmes, and the scientific underpinning of its horticulture. Bartram's Garden In 1728, John Bartram founded Bartram's Garden in Philadelphia, one of the continent's first botanical gardens. The garden is now managed as a historical site that includes a few original and many modern specimens as well as extensive archives and restored historical farm buildings. Plant classification The large number of plants needing description were often listed in garden catalogues; and at this time Carl Linnaeus established the system of binomial nomenclature which greatly facilitated the listing process. Names of plants were authenticated by dried plant specimens mounted on card (a or garden of dried plants) that were stored in buildings called herbaria, these taxonomic research institutions being frequently associated with the botanical gardens, many of which by then had "order beds" to display the classification systems being developed by botanists in the gardens' museums and herbaria. Botanical gardens had now become scientific collections, as botanists published their descriptions of the new exotic plants, and these were also recorded for posterity in detail by superb botanical illustrations. In this century, botanical gardens effectively dropped their medicinal function in favour of scientific and aesthetic priorities, and the term "botanic garden" came to be more closely associated with the herbarium, library (and later laboratories) housed there than with the living collectionson which little research was undertaken. 19th century The late 18th and early 19th centuries were marked by the establishment of tropical botanical gardens as a tool of colonial expansion (for trade and commerce and, secondarily, science) mainly by the British and Dutch, in India, South-east Asia and the Caribbean. This was also the time of Sir Joseph Banks's botanical collections during Captain James Cook's circumnavigations of the planet and his explorations of Oceania, which formed the last phase of plant introduction on a grand scale. Tropical botanical gardens There are currently about 230 tropical botanical gardens with a concentration in southern and south-eastern Asia. The first botanical garden founded in the tropics was the Pamplemousses Botanical Garden in Mauritius, established in 1735 to provide food for ships using the port, but later trialling and distributing many plants of economic importance. This was followed by the West Indies (Saint Vincent and the Grenadines Botanic Gardens, 1764) and in 1786 by the Acharya Jagadish Chandra Bose Botanical Garden in Calcutta, India founded during a period of prosperity when the city was a trading centre for the Dutch East India Company. Other gardens were constructed in Brazil (Rio de Janeiro Botanical Garden, 1808), Sri Lanka (Botanic Gardens of Peradeniya, 1821 and on a site dating back to 1371), Indonesia (Bogor Botanical Gardens, 1817 and Kebun Raya Cibodas, 1852), and Singapore (Singapore Botanical Gardens, 1822). These had a profound effect on the economy of the countries, especially in relation to the foods and medicines introduced. The importation of rubber trees to the Singapore Botanic Garden initiated the important rubber industry of the Malay Peninsula. At this time also, teak and tea were introduced to India and breadfruit, pepper and starfruit to the Caribbean. Included in the charter of these gardens was the investigation of the local flora for its economic potential to both the colonists and the local people. Many crop plants were introduced by or through these gardensoften in association with European botanical gardens such as Kew or Amsterdamand included cloves, tea, coffee, breadfruit, cinchona, sugar, cotton, palm oil and Theobroma cacao (for chocolate). During these times, the rubber plant was introduced to Singapore. Especially in the tropics, the larger gardens were frequently associated with a herbarium and museum of economy. The Botanical Garden of Peradeniya had considerable influence on the development of agriculture in Ceylon where the Para rubber tree () was introduced from Kew, which had itself imported the plant from South America. Other examples include cotton from the Chelsea Physic Garden to the Province of Georgia in 1732 and tea into India by Calcutta Botanic Garden. The transfer of germplasm between the temperate and tropical botanical gardens was undoubtedly responsible for the range of agricultural crops currently used in several regions of the tropics. Australia The first botanical gardens in Australia were founded early in the 19th century. The Royal Botanic Gardens, Sydney, 1816; the Royal Tasmanian Botanical Gardens, 1818; the Royal Botanic Gardens, Melbourne, 1845; Adelaide Botanic Gardens, 1854; and Brisbane Botanic Gardens, 1855. These were established essentially as colonial gardens of economic botany and acclimatisation. The Auburn Botanical Gardens, 1977, located in Sydney's western suburbs, are one of the popular and diverse botanical gardens in the Greater Western Sydney area. New Zealand Major botanical gardens in New Zealand include Dunedin Botanic Garden, 1863; Christchurch Botanic Gardens, 1863; Ōtari-Wilton's Bush, 1926; and Wellington Botanic Garden, 1868. Hong Kong Hong Kong Botanic Gardens, 1871 (renamed Hong Kong Zoological and Botanical Gardens in 1975), up from the Government Hill in Victoria City, Hong Kong Island. Japan The Koishikawa Botanical Garden in Tokyo, with its origin going back to the Tokugawa shogunate's ownership, became in 1877 part of the Tokyo Imperial University. Sri Lanka In Sri Lanka major botanical gardens include the Royal Botanic Gardens, Peradeniya (formally established in 1843), Hakgala Botanical Gardens (1861) and Henarathgoda Botanical Garden (1876). Ecuador Jardín Botánico de Quito is inside the Parque La Carolina is a 165.5-acre (670,000 m2) park in the centre of the Quito central business district, bordered by the avenues Río Amazonas, de los Shyris, Naciones Unidas, Eloy Alfaro, and de la República. The botanical garden of Quito is a park, a botanical garden, an arboretum and greenhouses of 18,600 square meters that is planned to increase, maintain the plants of the country (Ecuador is among the 17 richest countries in the world in the native species, a study on this matter). The Ecuadorian flora classified, determines the existence of 17,000 species) Egypt The Orman Garden, one of the most famous botanical gardens in Egypt, is located at Giza, in Cairo, and dates back to 1875. South Africa South Africa has ten national level botanical gardens all of which are overseen by the South African National Biodiversity Institute (SANBI). The oldest botanical garden in South Africa is the Durban Botanic Gardens which has been located on the same site since 1851. The Kirstenbosch National Botanical Garden is the most famous and developed garden in the country, established in 1913 on a site dating to 1848 and covering a 36 hectare area with an additional 528 hectares of mountainside wilderness that form part of the garden. Stellenbosch University Botanical Garden is the oldest university botanical garden in South Africa, and was established in 1922. Other botanical gardens in country include the Walter Sisulu National Botanical Garden, Harold Porter National Botanical Gardens and Karoo Desert National Botanical Garden. Some smaller gardens and parks that verge on being a botanical garden includes the Arderne Gardens in Cape Town founded in 1845. United States The first botanical garden in the United States, Bartram's Garden, was founded in 1730 near Philadelphia, and in the same year, the Linnaean Botanic Garden at Philadelphia itself. Presidents George Washington, Thomas Jefferson and James Madison, all experienced farmers, shared the dream of a national botanic garden for the collection, preservation and study of plants from around the world to contribute to the welfare of the American people paving the way for establishing the US Botanic Garden, right outside the nation's Capitol in Washington DC in 1820. In 1859, the Missouri Botanical Garden was founded at St Louis; it is now one of the world's leading gardens specializing in tropical plants. This was one of several popular American gardens, including Longwood Gardens (1798), Arnold Arboretum (1872), New York Botanical Garden (1891), Huntington Botanical Gardens (1906), Brooklyn Botanic Garden (1910), International Peace Garden (1932), and Fairchild Tropical Botanic Garden (1938). The first native plant garden in the United States was established in 1907 by Eloise Butler. US tax code provides for a substantial benefit to botanical gardens, this has led to a large number of entities declaring their campuses a botanical garden with little regard for veracity. Russia Russia has more gardens describing themselves as botanical gardens than any other country. Better-known gardens are Moscow University Botanic Garden ('the Apothecary Garden'), (1706), Saint Petersburg Botanical Garden, (1714); and Moscow Botanical Garden of Academy of Sciences, (1945). These gardens are notable for their structures that include sculptures, pavilions, bandstands, memorials, shadehouses, tea houses and such. Among the smaller gardens within Russia, one that is increasingly gaining prominence, is the Botanical Garden of Tver State University (1879) – the northernmost botanical Garden with an exhibition of steppe plants, only one of its kind in the Upper Volga. Ukraine Ukraine has about 30 botanical gardens. The most respected collections are Nikitsky Botanical Garden, Yalta, founded in 1812, M. M. Gryshko National Botanical Garden, a botanical garden of the National Academy of Sciences of Ukraine founded in 1936, and A. V. Fomin Botanical Garden of the Taras Shevchenko National University of Kyiv founded in 1839, which are located in Kyiv, the capital of Ukraine. 20th century Civic and municipal botanical gardens A large number of civic or municipal botanical gardens were founded in the 19th and 20th centuries. These did not develop scientific facilities or programmes, but the horticultural aspects were strong and the plants often labelled. They were botanical gardens in the sense of building up collections of plants and exchanging seeds with other gardens around the world, although their collection policies were determined by those in day-to-day charge of them. They tended to become little more than beautifully maintained parks and were, indeed, often under general parks administrations. Community engagement The second half of the 20th century saw increasingly sophisticated educational, visitor service, and interpretation services. Botanical gardens started to cater for many interests and their displays reflected this, often including botanical exhibits on themes of evolution, ecology or taxonomy, horticultural displays of attractive flowerbeds and herbaceous borders, plants from different parts of the world, special collections of plant groups such as bamboos or roses, and specialist glasshouse collections such as tropical plants, alpine plants, cacti and orchids, as well as the traditional herb gardens and medicinal plants. Specialised gardens like the Palmengarten in Frankfurt, Germany (1869), one of the world's leading orchid and succulent plant collections, have been very popular. There was a renewed interest in gardens of indigenous plants and areas dedicated to natural vegetation. With decreasing financial support from governments, revenue-raising public entertainment increased, including music, art exhibitions, special botanical exhibitions, theatre and film, this being supplemented by the advent of "Friends" organisations and the use of volunteer guides. Plant conservation Plant conservation and the heritage value of exceptional historic landscapes were treated with a growing sense of urgency. Specialist gardens were sometimes given a separate or adjoining site, to display native and indigenous plants. In the 1970s, gardens became focused on plant conservation. The Botanic Gardens Conservation Secretariat was established by the IUCN, and the World Conservation Union in 1987 with the aim of coordinating the plant conservation efforts of botanical gardens around the world. It maintains a database of rare and endangered species in botanical gardens' living collections. Many gardens hold ex situ conservation collections that preserve genetic variation. These may be held as seeds dried and stored at low temperature, or in tissue culture (such as the Kew Millennium Seedbank); as living plants, including those that are of special horticultural, historical or scientific interest (such as those in the National Plant Collection in the United Kingdom); or by managing and preserving areas of natural vegetation. Collections are often held and cultivated with the intention of reintroduction to their original habitats. The Center for Plant Conservation at St Louis, Missouri, coordinates the conservation of native North American species. Role and functions Many of the functions of botanical gardens have already been discussed in the sections above, which emphasise the scientific underpinning of botanical gardens with their focus on research, education and conservation. However, as multifaceted organisations, all sites have their own special interests. In a remarkable paper on the role of botanical gardens, Ferdinand von Mueller (1825–1896), the director of the Royal Botanic Gardens, Melbourne (1852–1873), stated, "in all cases the objects [of a botanical garden] must be mainly scientific and predominantly instructive". He then detailed many of the objectives being pursued by the world's botanical gardens in the middle of the 19th century, when European gardens were at their height. Many of these are listed below to give a sense of the scope of botanical gardens' activities at that time, and the ways in which they differed from parks or what he called "public pleasure gardens": availability of plants for scientific research display of plant diversity in form and use display of plants of particular regions (including local) plants sometimes grown within their particular families plants grown for their seed or rarity major timber (American English: lumber) trees plants of economic significance glasshouse plants of different climates all plants accurately labelled records kept of plants and their performance catalogues of holdings published periodically research facilities utilising the living collections studies in plant taxonomy examples of different vegetation types student education a herbarium selection and introduction of ornamental and other plants to commerce studies of plant chemistry (phytochemistry) report on the effects of plants on livestock at least one collector maintained doing field work Botanical gardens must find a compromise between the need for peace and seclusion, while at the same time satisfying the public need for information and visitor services that include restaurants, information centres and sales areas that bring with them rubbish, noise, and hyperactivity. Attractive landscaping and planting design sometimes compete with scientific interests — with science now often taking second place. Some gardens are now heritage landscapes that are subject to constant demand for new exhibits and exemplary environmental management. Many gardens now have plant shops selling flowers, herbs, and vegetable seedlings suitable for transplanting; many, like the UBC Botanical Garden and Centre for Plant Research and the Chicago Botanic Garden, have plant-breeding programs and introduce new plants to the horticultural trade. Future Botanical gardens are still being built, such as the first botanical garden in Oman, which will be one of the largest gardens in the world. Once completed, it will house the first large-scale cloud forest in a huge glasshouse. Development of botanical gardens in China over recent years has been remarkable, including the Hainan Botanical Garden of Tropical Economic Plants South China Botanical Garden at Guangzhou, the Xishuangbanna Botanical Garden of Tropical Plants and the Xiamen Botanic Garden, but in developed countries, many have closed for lack of financial support, this being especially true of botanical gardens attached to universities. The Palestine Museum of Natural History has a botanic garden, which has been described as a site of nation-building and resistance by Silvia Hassouna. Botanical gardens have always responded to the interests and values of the day. If a single function were to be chosen from the early literature on botanical gardens, it would be their scientific endeavour and, flowing from this, their instructional value. In their formative years, botanical gardens were gardens for physicians and botanists, but then they progressively became more associated with ornamental horticulture and the needs of the general public. The scientific reputation of a botanical garden is now judged by the publications coming out of herbaria and similar facilities, not by its living collections. The interest in economic plants now has less relevance, and the concern with plant classification systems has all but disappeared, while a fascination with the curious, beautiful and new seems unlikely to diminish. In recent times, the focus has been on creating an awareness of the threat to the Earth's ecosystems from human populations and its consequent need for biological and physical resources. Botanical gardens provide an excellent medium for communication between the world of botanical science and the general public. Education programs can help the public develop greater environmental awareness by understanding the meaning and importance of ideas like conservation and sustainability. Photo gallery Maps BGCI garden ID, Botanical gardens, Europe	Technology	Buildings and infrastructure	null
69442	https://en.wikipedia.org/wiki/Waterfall	Waterfall	A waterfall is any point in a river or stream where water flows over a vertical drop or a series of steep drops. Waterfalls also occur where meltwater drops over the edge of a tabular iceberg or ice shelf. Waterfalls can be formed in several ways, but the most common method of formation is that a river courses over a top layer of resistant bedrock before falling onto softer rock, which erodes faster, leading to an increasingly high fall. Waterfalls have been studied for their impact on species living in and around them. Humans have had a distinct relationship with waterfalls since prehistory, travelling to see them, exploring and naming them. They can present formidable barriers to navigation along rivers. Waterfalls are religious sites in many cultures. Since the 18th century, they have received increased attention as tourist destinations, sources of hydropower, andparticularly since the mid-20th centuryas subjects of research. Definition and terminology A waterfall is generally defined as a point in a river where water flows over a steep drop that is close to or directly vertical. In 2000 Mabin specified that "The horizontal distance between the positions of the lip and plunge pool should be no more than c 25% of the waterfall height." There are various types and methods to classify waterfalls. Some scholars have included rapids as a subsection. What actually constitutes a waterfall continues to be debated. Waterfalls are sometimes interchangeably referred to as "cascades" and "cataracts", though some sources specify a cataract as being a larger and more powerful waterfall and a cascade as being smaller. A plunge pool is a type of stream pool formed at the bottom of a waterfall. A waterfall may also be referred to as a "foss" or "force". Formation Waterfalls are commonly formed in the upper course of a river where lakes flow into valleys in steep mountains. A river sometimes flows over a large step in the rocks that may have been formed by a fault line. Waterfalls can occur along the edge of a glacial trough, where a stream or river flowing into a glacier continues to flow into a valley after the glacier has receded or melted. The large waterfalls in Yosemite Valley are examples of this phenomenon, which is referred to as a hanging valley. Another reason hanging valleys may form is where two rivers join and one is flowing faster than the other. When warm and cold water meets by a gorge in the ocean, large underwater waterfalls can form as the cold water rushes to the bottom. Caprock model The caprock model of waterfall formation states that the river courses over resistant bedrock, erosion happens slowly and is dominated by impacts of water-borne sediment on the rock, while downstream the erosion occurs more rapidly. As the watercourse increases its velocity at the edge of the waterfall, it may pluck material from the riverbed, if the bed is fractured or otherwise more erodible. Hydraulic jets and hydraulic jumps at the toe of a falls can generate large forces to erode the bed, especially when forces are amplified by water-borne sediment. Horseshoe-shaped falls focus the erosion to a central point, also enhancing riverbed change below a waterfall. A process known as "potholing" involves local erosion of a potentially deep hole in bedrock due to turbulent whirlpools spinning stones around on the bed, drilling it out. Sand and stones carried by the watercourse therefore increase erosion capacity. This causes the waterfall to carve deeper into the bed and to recede upstream. Often over time, the waterfall will recede back to form a canyon or gorge downstream as it recedes upstream, and it will carve deeper into the ridge above it. The rate of retreat for a waterfall can be as high as one-and-a-half metres per year. Often, the rock stratum just below the more resistant shelf will be of a softer type, meaning that undercutting due to splashback will occur here to form a shallow cave-like formation known as a rock shelter under and behind the waterfall. Eventually, the outcropping, more resistant cap rock will collapse under pressure to add blocks of rock to the base of the waterfall. These blocks of rock are then broken down into smaller boulders by attrition as they collide with each other, and they also erode the base of the waterfall by abrasion, creating a deep plunge pool in the gorge downstream. Streams can become wider and shallower just above waterfalls due to flowing over the rock shelf, and there is usually a deep area just below the waterfall because of the kinetic energy of the water hitting the bottom. However, a study of waterfalls systematics reported that waterfalls can be wider or narrower above or below a falls, so almost anything is possible given the right geological and hydrological setting. Waterfalls normally form in a rocky area due to erosion. After a long period of being fully formed, the water falling off the ledge will retreat, causing a horizontal pit parallel to the waterfall wall. Eventually, as the pit grows deeper, the waterfall collapses to be replaced by a steeply sloping stretch of river bed. In addition to gradual processes such as erosion, earth movement caused by earthquakes or landslides or volcanoes can lead to the formation of waterfalls. Ecology Waterfalls are an important factor in determining the distribution of lotic organisms such as fish and aquatic invertebrates, as they may restrict dispersal along streams. The presence or absence of certain species can have cascading ecological effects, and thus cause differences in trophic regimes above and below waterfalls. Certain aquatic plants and insects also specialize in the environment of the waterfall itself. A 2012 study of the Agbokim Waterfalls, has suggested that they hold biodiversity to a much higher extent than previously thought. Waterfalls also affect terrestrial species. They create a small microclimate in their immediate vicinity characterized by cooler temperatures and higher humidity than the surrounding region, which may support diverse communities of mosses and liverworts. Species of these plants may have disjunct populations at waterfall zones far from their core range. Waterfalls provide nesting cover for several species of bird, such as the black swift and white-throated dipper. These species preferentially nest in the space behind the falling water, which is thought to be a strategy to avoid predation. Types Some waterfalls are also distinct in that they do not flow continuously. Ephemeral waterfalls only flow after a rain or a significant snowmelt. Waterfalls can also be found underground and in oceans. Humans and waterfalls Research The geographer Andrew Goudie wrote in 2020 that waterfalls have received "surprisingly limited research." Alexander von Humboldt wrote about them in the 1820s. There is no name for the specific field of researching waterfalls, and in the published literature been described as "scattered", though it is popular to describe studying waterfalls as "waterfallology". An early paper written on waterfalls was published in 1884 by William Morris Davis, a geologist known as the "father of American geography". In the 1930s Edward Rashleigh published a pioneering work on waterfalls. In 1942 Oscar von Engeln wrote of the lack of research on waterfalls: As late as 1985 a scholar felt that "waterfalls remain a very much neglected aspect of river studies". Studies of waterfalls increased dramatically in the second half of the 20th century. Numerous waterfall guidebooks exist, and the World Waterfall Database is a website cataloging thousands of waterfalls. Exploration and naming Many explorers have visited waterfalls. European explorers recorded waterfalls they came across. In 1493, Christopher Columbus noted Carbet Falls in Guadeloupe, which was likely the first waterfall Europeans recorded seeing in the Americas. In the late 1600s, Louis Hennepin visited North America, providing early descriptions of Niagara Falls and the Saint Anthony Falls. The geographer Brian J. Hudson argues that it was uncommon to specifically name waterfalls until the 1700s. The trend of Europeans specifically naming falls was in tandem with increased scientific focus on nature, the rise of Romanticism, and increased importance of hydropower with the Industrial Revolution. European explorers often preferred to give waterfalls names in their own language; for instance, David Livingstone named Victoria Falls after Queen Victoria, though it was known by local peoples as Mosi-oa-Tunya. Many waterfalls have descriptive names which can come from the river they are on, places they are near, their features, or events that happened near them. Some countries that were colonized by European nations have taken steps to return names to waterfalls previously renamed by European explorers. Exploration of waterfalls continues; the Gocta Cataracts were first announced to the world in 2006. Waterfalls can pose major barriers to travel. Canals are sometimes built as a method to go around them, other times things must be physically carried around or a railway built. In 1885, the geographer George Chisholm wrote that, "The most signal example of the effect of waterfalls and rapids in retarding the development of civilisation is undoubtedly presented by the continent of Africa, the 'darkness' of which is almost entirely due to this cause." Development and tourism Waterfalls are often visited by people simply to see them. Hudson theorizes that they make good tourism sites because they are generally considered beautiful and are relatively uncommon. Activities at waterfalls can include bathing, swimming, photography, rafting, canyoning, abseiling, rock climbing, and ice climbing. Waterfalls can also be sites for generating hydroelectric power and can hold good fishing opportunities. Wealthy people were known to visit areas with features such as waterfalls at least as early as in Ancient Rome and China. However, many waterfalls were essentially inaccessible due to the treacherous terrain surrounding them until improvements began to be made such as paths to the falls, becoming common across the United Kingdom and America in the 1800s and continuing through the 1900s and into the 21st century. Remote waterfalls are now often visited by air travel. Human development has also threatened many waterfalls. For instance, the Guaíra Falls, once one of the most powerful waterfalls in the world, were submerged in 1982 by a human-made dam, as were the Ripon Falls in 1952. Conversely, other waterfalls have seen significantly lower water levels as a result of diversion for hydroelectricity, such as the Tyssestrengene in Norway. Development of the areas around falls as tourist attractions has also destroyed the natural scene around many of them. Waterfalls are included on thirty-eight World Heritage Sites and many others are protected by governments. In culture Waterfalls play a role in many cultures, as religious sites and subjects of art and music. Many artists have painted waterfalls and they are referenced in many songs, such as those of the Kaluli people in Papua New Guinea. Michael Harner titled his study of the Jivaroan peoples of Ecuador The Jivaro: People of the Sacred Waterfalls. Artists such as those of the Hudson River School and J. M. W. Turner and John Sell Cotman painted particularly notable pictures of waterfalls in the 19th century. One of the versions of the Shinto purification ceremony of misogi involves standing underneath a waterfall in ritual clothing. In Japan the Nachi Falls are a site of pilgrimage, as are falls near Tirupati, India, and the Saut-d'Eau, Haiti. The Otavalos use Piguchi waterfall as part of the Churru ritual which serves as a coming of age ceremony. Many waterfalls in Africa were places of worship for the native peoples and got their names from gods in the local religion. "In Chinese tradition, the waterfall represents" the season of autumn, yin, and the Chinese dragon's power over water that comes from the former two. List There are thousands of waterfalls in the world, though no exact number has been calculated. The World Waterfall Database lists 7,827 as of 2013, but this is likely incomplete; as noted by Hudson, over 90% of their listings are in North America. Many guidebooks to local waterfalls have been published. There is also no agreement how to measure the height of a waterfall, or even what constitutes one. Angel Falls in Venezuela is the tallest waterfall in the world, the Khone Phapheng Falls in Laos are the widest, and the Inga Falls on the Congo River are the biggest by flow rate, while the Dry Falls in Washington are the largest confirmed waterfalls ever. The highest known subterranean waterfall is in Vrtoglavica Cave in Slovenia. The Denmark Strait cataract is an undersea overflow which could be considered a "waterfall" under a very broad usage of that term; if so included, it is the largest known waterfall. Artificial waterfalls are water features or fountains that imitate a natural waterfall. The Cascata delle Marmore is the tallest artificially built waterfall at .	Physical sciences	Fluvial landforms	null
69453	https://en.wikipedia.org/wiki/Interstellar%20medium	Interstellar medium	The interstellar medium (ISM) is the matter and radiation that exists in the space between the star systems in a galaxy. This matter includes gas in ionic, atomic, and molecular form, as well as dust and cosmic rays. It fills interstellar space and blends smoothly into the surrounding intergalactic space. The energy that occupies the same volume, in the form of electromagnetic radiation, is the interstellar radiation field. Although the density of atoms in the ISM is usually far below that in the best laboratory vacuums, the mean free path between collisions is short compared to typical interstellar lengths, so on these scales the ISM behaves as a gas (more precisely, as a plasma: it is everywhere at least slightly ionized), responding to pressure forces, and not as a collection of non-interacting particles. The interstellar medium is composed of multiple phases distinguished by whether matter is ionic, atomic, or molecular, and the temperature and density of the matter. The interstellar medium is composed primarily of hydrogen, followed by helium with trace amounts of carbon, oxygen, and nitrogen. The thermal pressures of these phases are in rough equilibrium with one another. Magnetic fields and turbulent motions also provide pressure in the ISM, and are typically more important, dynamically, than the thermal pressure. In the interstellar medium, matter is primarily in molecular form and reaches number densities of 1012 molecules per m3 (1 trillion molecules per m3). In hot, diffuse regions, gas is highly ionized, and the density may be as low as 100 ions per m3. Compare this with a number density of roughly 1025 molecules per m3 for air at sea level, and 1016 molecules per m3 (10 quadrillion molecules per m3) for a laboratory high-vacuum chamber. Within our galaxy, by mass, 99% of the ISM is gas in any form, and 1% is dust. Of the gas in the ISM, by number 91% of atoms are hydrogen and 8.9% are helium, with 0.1% being atoms of elements heavier than hydrogen or helium, known as "metals" in astronomical parlance. By mass this amounts to 70% hydrogen, 28% helium, and 1.5% heavier elements. The hydrogen and helium are primarily a result of primordial nucleosynthesis, while the heavier elements in the ISM are mostly a result of enrichment (due to stellar nucleosynthesis) in the process of stellar evolution. The ISM plays a crucial role in astrophysics precisely because of its intermediate role between stellar and galactic scales. Stars form within the densest regions of the ISM, which ultimately contributes to molecular clouds and replenishes the ISM with matter and energy through planetary nebulae, stellar winds, and supernovae. This interplay between stars and the ISM helps determine the rate at which a galaxy depletes its gaseous content, and therefore its lifespan of active star formation. Voyager 1 reached the ISM on August 25, 2012, making it the first artificial object from Earth to do so. Interstellar plasma and dust will be studied until the estimated mission end date of 2025. Its twin Voyager 2 entered the ISM on November 5, 2018. Interstellar matter Table 1 shows a breakdown of the properties of the components of the ISM of the Milky Way. The three-phase model put forward the static two phase equilibrium model to explain the observed properties of the ISM. Their modeled ISM included a cold dense phase (T < 300 K), consisting of clouds of neutral and molecular hydrogen, and a warm intercloud phase (T ~ 104 K), consisting of rarefied neutral and ionized gas. added a dynamic third phase that represented the very hot (T ~ 106 K) gas that had been shock heated by supernovae and constituted most of the volume of the ISM. These phases are the temperatures where heating and cooling can reach a stable equilibrium. Their paper formed the basis for further study over the subsequent three decades. However, the relative proportions of the phases and their subdivisions are still not well understood. The basic physics behind these phases can be understood through the behaviour of hydrogen, since this is by far the largest constituent of the ISM. The different phases are roughly in pressure balance over most of the Galactic disk, since regions of excess pressure will expand and cool, and likewise under-pressure regions will be compressed and heated. Therefore, since P = n k T, hot regions (high T) generally have low particle number density n. Coronal gas has low enough density that collisions between particles are rare and so little radiation is produced, hence there is little loss of energy and the temperature can stay high for periods of hundreds of millions of years. In contrast, once the temperature falls to O(105 K) with correspondingly higher density, protons and electrons can recombine to form hydrogen atoms, emitting photons which take energy out of the gas, leading to runaway cooling. Left to itself this would produce the warm neutral medium. However, OB stars are so hot that some of their photons have energy greater than the Lyman limit, E > 13.6 eV, enough to ionize hydrogen. Such photons will be absorbed by, and ionize, any neutral hydrogen atom they encounter, setting up a dynamic equilibrium between ionization and recombination such that gas close enough to OB stars is almost entirely ionized, with temperature around 8000 K (unless already in the coronal phase), until the distance where all the ionizing photons are used up. This ionization front marks the boundary between the Warm ionized and Warm neutral medium. OB stars, and also cooler ones, produce many more photons with energies below the Lyman limit, which pass through the ionized region almost unabsorbed. Some of these have high enough energy (> 11.3 eV) to ionize carbon atoms, creating a C II ("ionized carbon") region outside the (hydrogen) ionization front. In dense regions this may also be limited in size by the availability of photons, but often such photons can penetrate throughout the neutral phase and only get absorbed in the outer layers of molecular clouds. Photons with E > 4 eV or so can break up molecules such as H2 and CO, creating a photodissociation region (PDR) which is more or less equivalent to the Warm neutral medium. These processes contribute to the heating of the WNM. The distinction between Warm and Cold neutral medium is again due to a range of temperature/density in which runaway cooling occurs. The densest molecular clouds have significantly higher pressure than the interstellar average, since they are bound together by their own gravity. When stars form in such clouds, especially OB stars, they convert the surrounding gas into the warm ionized phase, a temperature increase of several hundred. Initially the gas is still at molecular cloud densities, and so at vastly higher pressure than the ISM average: this is a classical H II region. The large overpressure causes the ionized gas to expand away from the remaining molecular gas (a Champagne flow), and the flow will continue until either the molecular cloud is fully evaporated or the OB stars reach the end of their lives, after a few millions years. At this point the OB stars explode as supernovas, creating blast waves in the warm gas that increase temperatures to the coronal phase (supernova remnants, SNR). These too expand and cool over several million years until they return to average ISM pressure. The ISM in different kinds of galaxy Most discussion of the ISM concerns spiral galaxies like the Milky Way, in which nearly all the mass in the ISM is confined to a relatively thin disk, typically with scale height about 100 parsecs (300 light years), which can be compared to a typical disk diameter of 30,000 parsecs. Gas and stars in the disk orbit the galactic centre with typical orbital speeds of 200 km/s. This is much faster than the random motions of atoms in the ISM, but since the orbital motion of the gas is coherent, the average motion does not directly affect structure in the ISM. The vertical scale height of the ISM is set in roughly the same way as the Earth's atmosphere, as a balance between the local gravitation field (dominated by the stars in the disk) and the pressure. Further from the disk plane, the ISM is mainly in the low-density warm and coronal phases, which extend at least several thousand parsecs away from the disk plane. This galactic halo or 'corona' also contains significant magnetic field and cosmic ray energy density. The rotation of galaxy disks influences ISM structures in several ways. Since the angular velocity declines with increasing distance from the centre, any ISM feature, such as giant molecular clouds or magnetic field lines, that extend across a range of radius are sheared by differential rotation, and so tend to become stretched out in the tangential direction; this tendency is opposed by interstellar turbulence (see below) which tends to randomize the structures. Spiral arms are due to perturbations in the disk orbits - essentially ripples in the disk, that cause orbits to alternately converge and diverge, compressing and then expanding the local ISM. The visible spiral arms are the regions of maximum density, and the compression often triggers star formation in molecular clouds, leading to an abundance of H II regions along the arms. Coriolis force also influences large ISM features. Irregular galaxies such as the Magellanic Clouds have similar interstellar mediums to spirals, but less organized. In elliptical galaxies the ISM is almost entirely in the coronal phase, since there is no coherent disk motion to support cold gas far from the center: instead, the scale height of the ISM must be comperable to the radius of the galaxy. This is consistent with the observation that there is little sign of current star formation in ellipticals. Some elliptical galaxies do show evidence for a small disk component, with ISM similar to spirals, buried close to their centers. The ISM of lenticular galaxies, as with their other properties, appear intermediate between spirals and ellipticals. Very close to the center of most galaxies (within a few hundred light years at most), the ISM is profoundly modified by the central supermassive black hole: see Galactic Center for the Milky Way, and Active galactic nucleus for extreme examples in other galaxies. The rest of this article will focus on the ISM in the disk plane of spirals, far from the galactic center. Structures Astronomers describe the ISM as turbulent, meaning that the gas has quasi-random motions coherent over a large range of spatial scales. Unlike normal turbulence, in which the fluid motions are highly subsonic, the bulk motions of the ISM are usually larger than the sound speed. Supersonic collisions between gas clouds cause shock waves which compress and heat the gas, increasing the sounds speed so that the flow is locally subsonic; thus supersonic turbulence has been described as 'a box of shocklets', and is inevitably associated with complex density and temperature structure. In the ISM this is further complicated by the magnetic field, which provides wave modes such as Alfvén waves which are often faster than pure sound waves: if turbulent speeds are supersonic but below the Alfvén wave speed, the behaviour is more like subsonic turbulence. Stars are born deep inside large complexes of molecular clouds, typically a few parsecs in size. During their lives and deaths, stars interact physically with the ISM. Stellar winds from young clusters of stars (often with giant or supergiant HII regions surrounding them) and shock waves created by supernovae inject enormous amounts of energy into their surroundings, which leads to hypersonic turbulence. The resultant structures – of varying sizes – can be observed, such as stellar wind bubbles and superbubbles of hot gas, seen by X-ray satellite telescopes or turbulent flows observed in radio telescope maps. Stars and planets, once formed, are unaffected by pressure forces in the ISM, and so do not take part in the turbulent motions, although stars formed in molecular clouds in a galactic disk share their general orbital motion around the galaxy center. Thus stars are usually in motion relative to their surrounding ISM. The Sun is currently traveling through the Local Interstellar Cloud, an irregular clump of the warm neutral medium a few parsecs across, within the low-density Local Bubble, a 100-parsec radius region of coronal gas. In October 2020, astronomers reported a significant unexpected increase in density in the space beyond the Solar System as detected by the Voyager 1 and Voyager 2 space probes. According to the researchers, this implies that "the density gradient is a large-scale feature of the VLISM (very local interstellar medium) in the general direction of the heliospheric nose". Interaction with interplanetary medium The interstellar medium begins where the interplanetary medium of the Solar System ends. The solar wind slows to subsonic velocities at the termination shock, 90–100 astronomical units from the Sun. In the region beyond the termination shock, called the heliosheath, interstellar matter interacts with the solar wind. Voyager 1, the farthest human-made object from the Earth (after 1998), crossed the termination shock December 16, 2004 and later entered interstellar space when it crossed the heliopause on August 25, 2012, providing the first direct probe of conditions in the ISM . Interstellar extinction Dust grains in the ISM are responsible for extinction and reddening, the decreasing light intensity and shift in the dominant observable wavelengths of light from a star. These effects are caused by scattering and absorption of photons and allow the ISM to be observed with the naked eye in a dark sky. The apparent rifts that can be seen in the band of the Milky Way – a uniform disk of stars – are caused by absorption of background starlight by dust in molecular clouds within a few thousand light years from Earth. This effect decreases rapidly with increasing wavelength ("reddening" is caused by greater absorption of blue than red light), and becomes almost negligible at mid-infrared wavelengths (> 5 μm). Extinction provides one of the best ways of mapping the three-dimensional structure of the ISM, especially since the advent of accurate distances to millions of stars from the Gaia mission. The total amount of dust in front of each star is determined from its reddening, and the dust is then located along the line of sight by comparing the dust column density in front of stars projected close together on the sky, but at different distances. By 2022 it was possible to generate a map of ISM structures within 3 kpc (10,000 light years) of the Sun. Far ultraviolet light is absorbed effectively by the neutral hydrogen gas in the ISM. Specifically, atomic hydrogen absorbs very strongly at about 121.5 nanometers, the Lyman-alpha transition, and also at the other Lyman series lines. Therefore, it is nearly impossible to see light emitted at those wavelengths from a star farther than a few hundred light years from Earth, because most of it is absorbed during the trip to Earth by intervening neutral hydrogen. All photons with wavelength < 91.6 nm, the Lyman limit, can ionize hydrogen and are also very strongly absorbed. The absorption gradually decreases with increasing photon energy, and the ISM begins to become transparent again in soft X-rays, with wavelengths shorter than about 1 nm. Heating and cooling The ISM is usually far from thermodynamic equilibrium. Collisions establish a Maxwell–Boltzmann distribution of velocities, and the 'temperature' normally used to describe interstellar gas is the 'kinetic temperature', which describes the temperature at which the particles would have the observed Maxwell–Boltzmann velocity distribution in thermodynamic equilibrium. However, the interstellar radiation field is typically much weaker than a medium in thermodynamic equilibrium; it is most often roughly that of an A star (surface temperature of ~10,000 K) highly diluted. Therefore, bound levels within an atom or molecule in the ISM are rarely populated according to the Boltzmann formula . Depending on the temperature, density, and ionization state of a portion of the ISM, different heating and cooling mechanisms determine the temperature of the gas. Heating mechanisms Heating by low-energy cosmic rays The first mechanism proposed for heating the ISM was heating by low-energy cosmic rays. Cosmic rays are an efficient heating source able to penetrate in the depths of molecular clouds. Cosmic rays transfer energy to gas through both ionization and excitation and to free electrons through Coulomb interactions. Low-energy cosmic rays (a few MeV) are more important because they are far more numerous than high-energy cosmic rays. Photoelectric heating by grains The ultraviolet radiation emitted by hot stars can remove electrons from dust grains. The photon is absorbed by the dust grain, and some of its energy is used to overcome the potential energy barrier and remove the electron from the grain. This potential barrier is due to the binding energy of the electron (the work function) and the charge of the grain. The remainder of the photon's energy gives the ejected electron kinetic energy which heats the gas through collisions with other particles. A typical size distribution of dust grains is n(r) ∝ r, where r is the radius of the dust particle. Assuming this, the projected grain surface area distribution is πrn(r) ∝ r. This indicates that the smallest dust grains dominate this method of heating. Photoionization When an electron is freed from an atom (typically from absorption of a UV photon) it carries kinetic energy away of the order E − E. This heating mechanism dominates in H II regions, but is negligible in the diffuse ISM due to the relative lack of neutral carbon atoms. X-ray heating X-rays remove electrons from atoms and ions, and those photoelectrons can provoke secondary ionizations. As the intensity is often low, this heating is only efficient in warm, less dense atomic medium (as the column density is small). For example, in molecular clouds only hard x-rays can penetrate and x-ray heating can be ignored. This is assuming the region is not near an x-ray source such as a supernova remnant. Chemical heating Molecular hydrogen (H2) can be formed on the surface of dust grains when two H atoms (which can travel over the grain) meet. This process yields 4.48 eV of energy distributed over the rotational and vibrational modes, kinetic energy of the H2 molecule, as well as heating the dust grain. This kinetic energy, as well as the energy transferred from de-excitation of the hydrogen molecule through collisions, heats the gas. Grain-gas heating Collisions at high densities between gas atoms and molecules with dust grains can transfer thermal energy. This is not important in HII regions because UV radiation is more important. It is also less important in diffuse ionized medium due to the low density. In the neutral diffuse medium grains are always colder, but do not effectively cool the gas due to the low densities. Grain heating by thermal exchange is very important in supernova remnants where densities and temperatures are very high. Gas heating via grain-gas collisions is dominant deep in giant molecular clouds (especially at high densities). Far infrared radiation penetrates deeply due to the low optical depth. Dust grains are heated via this radiation and can transfer thermal energy during collisions with the gas. A measure of efficiency in the heating is given by the accommodation coefficient: where T is the gas temperature, Td the dust temperature, and T2 the post-collision temperature of the gas atom or molecule. This coefficient was measured by as α = 0.35. Other heating mechanisms A variety of macroscopic heating mechanisms are present including: Gravitational collapse of a cloud Supernova explosions Stellar winds Expansion of H II regions Magnetohydrodynamic waves created by supernova remnants Cooling mechanisms Fine structure cooling The process of fine structure cooling is dominant in most regions of the Interstellar Medium, except regions of hot gas and regions deep in molecular clouds. It occurs most efficiently with abundant atoms having fine structure levels close to the fundamental level such as: C II and O I in the neutral medium and O II, O III, N II, N III, Ne II and Ne III in H II regions. Collisions will excite these atoms to higher levels, and they will eventually de-excite through photon emission, which will carry the energy out of the region. Cooling by permitted lines At lower temperatures, more levels than fine structure levels can be populated via collisions. For example, collisional excitation of the n = 2 level of hydrogen will release a Ly-α photon upon de-excitation. In molecular clouds, excitation of rotational lines of CO is important. Once a molecule is excited, it eventually returns to a lower energy state, emitting a photon which can leave the region, cooling the cloud. Observations of the ISM Despite its extremely low density, photons generated in the ISM are prominent in nearly all bands of the electromagnetic spectrum. In fact the optical band, on which astronomers relied until well into the 20th century, is the one in which the ISM is least obvious. Ionized gas radiates at a broad range of energies via bremsstrahlung. For gas in the warm phase (104 K) this is mostly detected in microwaves, while bremsstrahlung from the million-kelvin coronal gas is prominent in soft X-rays. In addition, many spectral lines are produced, including the ones significant for cooling mentioned in the previous section. One of these, a forbidden line of doubly-ionized oxygen, gives many nebulae their apparent green colour in visual observations, and was once thought to be a new element, nebulium. Spectral lines from highly excited states of hydrogen are detectable at infra-red and longer wavelengths, down to radio recombination lines which, unlike optical lines, are not absorbed by dust and so can trace ionized regions throughout the disk of the Galaxy. Coronal gas emits a different set of lines, since atoms are stripped of a larger fraction of their electrons at its high temperature. The warm neutral medium produces most of the 21-cm line emission from hydrogen detected by radio telescopes, although atomic hydrogen in the cold neutral medium also contributes, both in emission and by absorption of photons from background warm gas ('H I self-absorption', HISA). While not important for cooling, the 21-cm line is easily observable at high spectral and angular resolution, giving us our most detailed view of the WNM. Molecular clouds are detected via spectral lines produced by changes in the rotational quantum state of small molecules, especially carbon monoxide, CO. The most widely used line is at 115 GHz, corresponding to the change from 1 to 0 quanta of angular momentum. Hundreds of other molecules have been detected, each with many lines, which allows physical and chemical processes in molecular clouds to be traced in some detail. These lines are most common at millimetre and sub-mm wavelengths. By far the most common molecule in molecular clouds, H2, is usually not directly observable, as it stays in its ground state except when excited by rare events such as interstellar shock waves. There is some 'dark gas', regions where hydrogen is in molecular form and therefore does not emit the 21-cm line, but CO molecules are broken up so the CO lines are also not present. These regions are inferred from the presence of dust grains with no matching line emission from gas. Interstellar dust grains re-emit the energy they absorb from starlight as quasi-blackbody emission in the far infrared, corresponding to typical dust grain temperatures of 20–100 K. Very small grains, essentially fragments of graphene bonded to hydrogen atoms around their edges (polycyclic aromatic hydrocarbons, PAHs), emit numerous spectral lines in the mid-infrared, at wavelengths around 10 micron. Nanometre-sized grains can be spun up to rotate at GHz frequencies by a collision with a single ultraviolet photon, and dipole radiation from such spinning grains is believed to be the source of anomalous microwave emission. Cosmic rays generate gamma-ray photons when they collide with atomic nuclei in ISM clouds. The electrons amongst cosmic ray particles collide with a small fraction of photons in the interstellar radiation field and the cosmic microwave background and bump up the photon energies to X-rays and gamma-rays, via inverse Compton scattering. Due to the galactic magnetic field, charged particles follow spiral paths, and for cosmic-ray electrons this spiralling motion generates synchrotron radiation which is very bright at low radio frequencies. Radiowave propagation Radio waves are affected by the plasma properties of the ISM. The lowest frequency radio waves, below ≈ 0.1 MHz, cannot propagate through the ISM since they are below its plasma frequency. At higher frequencies, the plasma has a significant refractive index, decreasing with increasing frequency, and also dependent on the density of free electrons. Random variations in the electron density cause interstellar scintillation, which broadens the apparent size of distant radio sources seen through the ISM, with the broadening decreasing with frequency squared. The variation of refractive index with frequency causes the arrival times of pulses from pulsars and Fast radio bursts to be delayed at lower frequencies (dispersion). The amount of delay is proportional to the column density of free electrons (Dispersion measure, DM), which is useful for both mapping the distribution of ionized gas in the Galaxy and estimating distances to pulsars (more distant ones have larger DM). A second propagation effect is Faraday rotation, which affects linearly polarized radio waves, such as those produced by synchrotron radiation, one of the most common sources of radio emission in astrophysics. Faraday rotation depends on both the electron density and the magnetic field strength, and so is used as a probe of the interstellar magnetic field. The ISM is generally very transparent to radio waves, allowing unimpeded observations right through the disk of the Galaxy. There are a few exceptions to this rule. The most intense spectral lines in the radio spectrum can become opaque, so that only the surface of the line-emitting cloud is visible. This mainly affects the carbon monoxide lines at millimetre wavelengths that are used to trace molecular clouds, but the 21-cm line from neutral hydrogen can become opaque in the cold neutral medium. Such absorption only affects photons at the line frequencies: the clouds are otherwise transparent. The other significant absorption process occurs in dense ionized regions. These emit photons, including radio waves, via thermal bremsstrahlung. At short wavelengths, typically microwaves, these are quite transparent, but their brightness approaches the black body limit as , and at wavelengths long enough that this limit is reached, they become opaque. Thus metre-wavelength observations show H II regions as cool spots blocking the bright background emission from Galactic synchrotron radiation, while at decametres the entire galactic plane is absorbed, and the longest radio waves observed, 1 km, can only propagate 10-50 parsecs through the Local Bubble. The frequency at which a particular nebula becomes optically thick depends on its emission measure , the column density of squared electron number density. Exceptionally dense nebulae can become optically thick at centimetre wavelengths: these are just-formed and so both rare and small ('Ultra-compact H II regions') The general transparency of the ISM to radio waves, especially microwaves, may seem surprising since radio waves at frequencies > 10 GHz are significantly attenuated by Earth's atmosphere (as seen in the figure). But the column density through the atmosphere is vastly larger than the column through the entire Galaxy, due to the extremely low density of the ISM. History of knowledge of interstellar space The word 'interstellar' (between the stars) was coined by Francis Bacon in the context of the ancient theory of a literal sphere of fixed stars. Later in the 17th century, when the idea that stars were scattered through infinite space became popular, it was debated whether that space was a true vacuum or filled with a hypothetical fluid, sometimes called aether, as in René Descartes' vortex theory of planetary motions. While vortex theory did not survive the success of Newtonian physics, an invisible luminiferous aether was re-introduced in the early 19th century as the medium to carry light waves; e.g., in 1862 a journalist wrote: "this efflux occasions a thrill, or vibratory motion, in the ether which fills the interstellar spaces." In 1864, William Huggins used spectroscopy to determine that a nebula is made of gas. Huggins had a private observatory with an 8-inch telescope, with a lens by Alvan Clark; but it was equipped for spectroscopy, which enabled breakthrough observations. From around 1889, Edward Barnard pioneered deep photography of the sky, finding many 'holes in the Milky Way'. At first he compared them to sunspots, but by 1899 was prepared to write: "One can scarcely conceive a vacancy with holes in it, unless there is nebulous matter covering these apparently vacant places in which holes might occur". These holes are now known as dark nebulae, dusty molecular clouds silhouetted against the background star field of the galaxy; the most prominent are listed in his Barnard Catalogue. The first direct detection of cold diffuse matter in interstellar space came in 1904, when Johannes Hartmann observed the binary star Mintaka (Delta Orionis) with the Potsdam Great Refractor. Hartmann reported that absorption from the "K" line of calcium appeared "extraordinarily weak, but almost perfectly sharp" and also reported the "quite surprising result that the calcium line at 393.4 nanometres does not share in the periodic displacements of the lines caused by the orbital motion of the spectroscopic binary star". The stationary nature of the line led Hartmann to conclude that the gas responsible for the absorption was not present in the atmosphere of the star, but was instead located within an isolated cloud of matter residing somewhere along the line of sight to this star. This discovery launched the study of the interstellar medium. Interstellar gas was further confirmed by Slipher in 1909, and then by 1912 interstellar dust was confirmed by Slipher. Interstellar sodium was detected by Mary Lea Heger in 1919 through the observation of stationary absorption from the atom's "D" lines at 589.0 and 589.6 nanometres towards Delta Orionis and Beta Scorpii. In the series of investigations, Viktor Ambartsumian introduced the now commonly accepted notion that interstellar matter occurs in the form of clouds. Subsequent observations of the "H" and "K" lines of calcium by revealed double and asymmetric profiles in the spectra of Epsilon and Zeta Orionis. These were the first steps in the study of the very complex interstellar sightline towards Orion. Asymmetric absorption line profiles are the result of the superposition of multiple absorption lines, each corresponding to the same atomic transition (for example the "K" line of calcium), but occurring in interstellar clouds with different radial velocities. Because each cloud has a different velocity (either towards or away from the observer/Earth), the absorption lines occurring within each cloud are either blue-shifted or red-shifted (respectively) from the lines' rest wavelength through the Doppler Effect. These observations confirming that matter is not distributed homogeneously were the first evidence of multiple discrete clouds within the ISM. The growing evidence for interstellar material led to comment: "While the interstellar absorbing medium may be simply the ether, yet the character of its selective absorption, as indicated by Kapteyn, is characteristic of a gas, and free gaseous molecules are certainly there, since they are probably constantly being expelled by the Sun and stars." The same year, Victor Hess's discovery of cosmic rays, highly energetic charged particles that rain onto the Earth from space, led others to speculate whether they also pervaded interstellar space. The following year, the Norwegian explorer and physicist Kristian Birkeland wrote: "It seems to be a natural consequence of our points of view to assume that the whole of space is filled with electrons and flying electric ions of all kinds. We have assumed that each stellar system in evolutions throws off electric corpuscles into space. It does not seem unreasonable therefore to think that the greater part of the material masses in the universe is found, not in the solar systems or nebulae, but in 'empty' space" . noted that "it could scarcely have been believed that the enormous gaps between the stars are completely void. Terrestrial aurorae are not improbably excited by charged particles emitted by the Sun. If the millions of other stars are also ejecting ions, as is undoubtedly true, no absolute vacuum can exist within the galaxy." In September 2012, NASA scientists reported that polycyclic aromatic hydrocarbons (PAHs), subjected to interstellar medium (ISM) conditions, are transformed, through hydrogenation, oxygenation and hydroxylation, to more complex organics, "a step along the path toward amino acids and nucleotides, the raw materials of proteins and DNA, respectively". Further, as a result of these transformations, the PAHs lose their spectroscopic signature, which could be one of the reasons "for the lack of PAH detection in interstellar ice grains, particularly the outer regions of cold, dense clouds or the upper molecular layers of protoplanetary disks." In February 2014, NASA announced a greatly upgraded database for tracking polycyclic aromatic hydrocarbons (PAHs) in the universe. According to scientists, more than 20% of the carbon in the universe may be associated with PAHs, possible starting materials for the formation of life. PAHs seem to have been formed shortly after the Big Bang, are widespread throughout the universe, and are associated with new stars and exoplanets. In April 2019, scientists, working with the Hubble Space Telescope, reported the confirmed detection of the large and complex ionized molecules of buckminsterfullerene (C60) (also known as "buckyballs") in the interstellar medium spaces between the stars. In September 2020, evidence was presented of solid-state water in the interstellar medium, and particularly, of water ice mixed with silicate grains in cosmic dust grains.	Physical sciences	Basics_3	null
69562	https://en.wikipedia.org/wiki/Vasectomy	Vasectomy	Vasectomy is an elective surgical procedure that results in male sterilization, often as a means of permanent contraception. During the procedure, the male vasa deferentia are cut and tied or sealed so as to prevent sperm from entering into the urethra and thereby prevent fertilization of a female through sexual intercourse. Vasectomies are usually performed in a physician's office, medical clinic, or, when performed on a non-human animal, in a veterinary clinic. Hospitalization is not normally required as the procedure is not complicated, the incisions are small, and the necessary equipment routine. There are several methods by which a surgeon might complete a vasectomy procedure, all of which occlude (i.e., "seal") at least one side of each vas deferens. To help reduce anxiety and increase patient comfort, those who have an aversion to needles may consider a "no-needle" application of anesthesia while the 'no-scalpel' or 'open-ended' techniques help to accelerate recovery times and increase the chance of healthy recovery. Due to the simplicity of the surgery, a vasectomy usually takes less than 30 minutes to complete. After a short recovery at the doctor's office (usually less than an hour), the patient is sent home to rest. Because the procedure is minimally invasive, many vasectomy patients find that they can resume their typical sexual behavior within a week, and do so with little or no discomfort. Because the procedure is considered a permanent method of contraception and is not easily reversed, patients are frequently counseled and advised to consider how the long-term outcome of a vasectomy might affect them both emotionally and physically. The procedure is not typically encouraged for young single childless people as their risk of later regret is higher as chances of biological parenthood are thereby permanently reduced, often completely. A vasectomy without the patient's consent or knowledge is considered forced sterilization. Medical uses A vasectomy is done to prevent fertility in males. It ensures that in most cases the person will be sterile after confirmation of success following surgery. The procedure is regarded as permanent because vasectomy reversal is costly and often does not restore the male's sperm count or sperm motility to prevasectomy levels. Those with vasectomies have a very small (nearly zero) chance of successfully impregnating someone, but a vasectomy does not protect against sexually transmitted infections (STIs). After vasectomy, the testes remain in the scrotum where Leydig cells continue to produce testosterone and other male hormones that continue to be secreted into the bloodstream. When the vasectomy is complete, sperm cannot exit the body through the penis. Sperm is still produced by the testicles but is broken down and absorbed by the body. Much fluid content is absorbed by membranes in the epididymis, and much solid content is broken down by the responding macrophages and reabsorbed via the bloodstream. After vasectomy, the membranes must increase in size to absorb and store more fluid; this triggering of the immune system causes more macrophages to be recruited to break down and reabsorb more solid content. Within one year after vasectomy, sixty to seventy percent of those vasectomized develop antisperm antibodies. In some cases, vasitis nodosa, a benign proliferation of the ductular epithelium, can also result. The accumulation of sperm increases pressure in the vas deferens and epididymis. The entry of the sperm into the scrotum can cause sperm granulomas to be formed by the body to contain and absorb the sperm which the body will treat as a foreign biological substance (much like a virus or bacterium). Efficacy Vasectomy is the most effective permanent form of contraception available to males. (Removing the entire vas deferens would very likely be more effective, but it is not something that is regularly done.) In nearly every way that vasectomy can be compared to tubal ligation it has a more positive outlook. Vasectomy is more cost effective, less invasive, has techniques that are emerging that may facilitate easier reversal, and has a much lower risk of postoperative complications. Early failure rates, i.e. pregnancy within a few months after vasectomy, typically result from unprotected sexual intercourse too soon after the procedure while some sperm continue to pass through the vasa deferentia. Most physicians and surgeons who perform vasectomies recommend one (sometimes two) postprocedural semen specimens to verify a successful vasectomy; however, many people fail to return for verification tests citing inconvenience, embarrassment, forgetfulness, or certainty of sterility. In January 2008, the FDA cleared a home test called SpermCheck Vasectomy that allows patients to perform postvasectomy confirmation tests themselves; however, compliance for postvasectomy semen analysis in general remains low. Late failure, i.e. pregnancy following spontaneous recanalization of the vasa deferentia, has also been documented. This occurs because the epithelium of the vas deferens (similar to the epithelium of some other human body parts) is capable of regenerating and creating a new tube if the vas deferens is damaged and/or severed. Even when as much as five centimeters (or two inches) of the vas deferens is removed, the vas deferens can still grow back together and become reattached—thus allowing sperm to once again pass and flow through the vas deferens, restoring one's fertility. The Royal College of Obstetricians and Gynaecologists states there is a generally agreed-upon rate of late failure of about one in 2000 vasectomies—better than tubal ligations for which the failure rate is one in every 200 to 300 cases. A 2005 review including both early and late failures described a total of 183 recanalizations from 43,642 vasectomies (0.4%), and 60 pregnancies after 92,184 vasectomies (0.07%). Complications and patient concerns Short-term possible complications include infection, bruising and bleeding into the scrotum resulting in a collection of blood known as a hematoma. A study in 2012 demonstrated an infection rate of 2.5% postvasectomy. The stitches on the small incisions required are prone to irritation, though this can be minimized by covering them with gauze or small adhesive bandages. The primary long-term complications are chronic pain conditions or syndromes that can affect any of the scrotal, pelvic or lower-abdominal regions, collectively known as post-vasectomy pain syndrome. Though vasectomy results in increases in circulating immune complexes, these increases are transient. Data based on animal and human studies indicate these changes do not result in increased incidence of atherosclerosis. Complications not withstanding, many men express concerns regarding potential adverse effects of vasectomy, including Cancer. The risk of testicular cancer is not affected by vasectomy. In 2014 the AUA reaffirmed that vasectomy is not a risk factor for prostate cancer and that it is not necessary for physicians to routinely discuss prostate cancer in their preoperative counseling of vasectomy patients. There remains ongoing debate regarding whether vasectomy is associated with prostate cancer. A 2017 meta-analysis found no statistically significant increase in risk. A 2019 study of 2.1 million Danish males found that vasectomy increased their incidence of prostate cancer by 15%. A 2020 meta-analysis found that vasectomy increased the incidence by 9%. Other recent studies agree on the 15% increase in risk of developing prostate cancer, but found that people who get a vasectomy are not more likely to die from prostate cancer than those without a vasectomy. A vasectomy will not Impact sexual performance: A vasectomy does not affect libido or masculinity aside from its contraceptive effect. Some men have even reported increased sexual satisfaction post-vasectomy. It won't cause permanent damage to sexual organs: The risk of injury to the testicles, penis, or other reproductive organs during surgery is minimal. Although extremely rare, damage to the blood supply could potentially lead to testicular loss, but this occurrence is unlikely with a skilled surgeon. It won't increase the risk of heart disease: There is no established connection between vasectomy and heart-related issues. Vasectomy will not result in severe pain: While minor discomfort such as pulling or tugging sensations may occur during the procedure, severe pain is uncommon. Post-surgery, most men experience minor pain that typically resolves within a few days. In rare cases, some men report chronic post surgery pain, Post Vasectomy Pain Syndrome. Postvasectomy pain Post-vasectomy pain syndrome is a chronic and sometimes debilitating condition that may develop immediately or several years after vasectomy. The most robust study of post-vasectomy pain, according to the American Urology Association's Vasectomy Guidelines 2012 (amended 2015) surveyed people just before their vasectomy and again seven months later. Of those that responded and who said they did not have any scrotal pain prior to vasectomy, 7% had scrotal pain seven months later which they described as "Mild, a bit of a nuisance", 1.6% had pain that was "Moderate, require painkillers" and 0.9% had pain that was "quite severe and noticeably affecting their quality of life". Post-vasectomy pain can be constant orchialgia or epididymal pain (epididymitis), or it can be pain that occurs only at particular times such as with sexual intercourse, ejaculation, or physical exertion. Psychological effects Approximately 90% are generally reported in reviews as being satisfied with having had a vasectomy, while 7–10% of people regret their decision. For those in relationships, regret was less common when both people in the relationship agreed on the procedure. Younger people who receive a vasectomy are significantly more likely to regret and seek a reversal of their vasectomy, with one study showing people in their twenties being 12.5 times more likely to undergo a vasectomy reversal later in life (and including some who chose sterilization at a young age). Pre-vasectomy counseling is often emphasised for younger patients. Dementia An association between vasectomy and primary progressive aphasia, a rare variety of frontotemporal dementia, was reported. However, it is doubtful that there is a causal relationship. The putative mechanism is a cross-reactivity between brain and sperm, including the shared presence of neural surface antigens. In addition, the cytoskeletal tau protein has been found only to exist outside of the CNS in the manchette of sperm. Procedure The traditional incision approach of vasectomy involves numbing of the scrotum with local anesthetic (although some people's physiology may make access to the vas deferens more difficult in which case general anesthesia may be recommended) after which a scalpel is used to make two small incisions, one on each side of the scrotum at a location that allows the surgeon to bring each vas deferens to the surface for excision. The vasa deferentia are cut (sometimes a section may be removed altogether), separated, and then at least one side is sealed by ligating (suturing), cauterizing (electrocauterization), or clamping. There are several variations to this method that may improve healing, effectiveness, and which help mitigate long-term pain such as post-vasectomy pain syndrome or epididymitis, however the data supporting one over another are limited. Fascial interposition: Recanalization of the vas deferens is a known cause of vasectomy failure(s). Fascial interposition ("FI"), in which a tissue barrier is placed between the cut ends of the vas by suturing, may help to prevent this type of failure, increasing the overall success rate of vasectomy while leaving the testicular end within the confines of the fascia. The fascia is a fibrous protective sheath that surrounds the vas deferens as well as all other body muscle tissue. This method, when combined with intraluminal cautery (where one or both sides of the vas deferens are electrically "burned" closed to prevent recanalization), has been shown to increase the success rate of vasectomy procedures. No-needle anesthesia: Fear of needles for injection of local anesthesia is well known. In 2005, a method of local anesthesia was introduced for vasectomy which allows the surgeon to apply it painlessly with a special jet-injection tool, as opposed to traditional needle application. The numbing agent is forced/pushed onto and deep enough into the scrotal tissue to allow for a virtually pain-free surgery. Lidocaine applied in this manner typically achieves anesthesia within 10 to 20 seconds. Initial surveys show a very high satisfaction rate amongst vasectomy patients. Once the effects of no-needle anesthesia set in, the vasectomy procedure is performed in the routine manner. However, unlike in conventional local anesthesia where needles and syringes are used on one patient only, the applicator is not single use and can only be properly disinfected by autoclaving. No-scalpel vasectomy (NSV): Also known as a "key-hole" vasectomy, is a vasectomy in which a sharp hemostat (as opposed to a scalpel) is used to puncture the scrotum. This method has come into widespread use as the resulting smaller "incision" or puncture wound typically limits bleeding and hematomas. Also the smaller wound has less chance of infection, resulting in faster healing times compared to the larger/longer incisions made with a scalpel. The surgical wound created by the no-scalpel method usually does not require stitches. NSV is the most commonly performed type of minimally invasive vasectomy, and both describe the method of vasectomy that leads to access of the vas deferens. Open-ended vasectomy: In this procedure the testicular end of the vas deferens is not sealed, which allows continued streaming of sperm into the scrotum. This method may avoid testicular pain resulting from increased back-pressure in the epididymis. Studies suggest that this method may reduce long-term complications such as post-vasectomy pain syndrome. Vas irrigation: Injections of sterile water or euflavine (which kills sperm) are put into the distal portion of the vas at the time of surgery which then brings about a near-immediate sterile ("azoospermatic") condition. The use of euflavine does however, tend to decrease time (or, number of ejaculations) to azoospermia vs. the water irrigation by itself. This additional step in the vasectomy procedure, (and similarly, fascial interposition), has shown positive results but is not as prominently in use, and few surgeons offer it as part of their vasectomy procedure. Other techniques The following vasectomy methods have purportedly had a better chance of later reversal but have seen less use by virtue of known higher failure rates (i.e., recanalization). An earlier clip device, the VasClip, is no longer on the market, due to unacceptably high failure rates. The VasClip method, though considered reversible, has had a higher cost and resulted in lower success rates. Also, because the vasa deferentia are not cut or tied with this method, it could technically be classified as other than a vasectomy. Vasectomy reversal (and the success thereof) was conjectured to be higher as it only required removing the Vas-Clip device. This method achieved limited use, and scant reversal data are available. Vas occlusion techniques Injected plugs: There are two types of injected plugs which can be used to block the vasa deferentia. Medical-grade polyurethane (MPU) or medical-grade silicone rubber (MSR) starts as a liquid polymer that is injected into the vas deferens after which the liquid is clamped in place until is solidifies (usually in a few minutes). Intra-vas device: The vasa deferentia can also be occluded by an intra-vas device (IVD). A small cut is made in the lower abdomen after which a soft silicone or urethane plug is inserted into each vas tube thereby blocking (occluding) sperm. This method allows for the vas to remain intact. IVD technique is done in an out-patient setting with local anesthetic, similar to a traditional vasectomy. IVD reversal can be performed under the same conditions making it much less costly than vasovasostomy which can require general anesthesia and longer surgery time. Both vas occlusion techniques require the same basic patient setup: local anesthesia, puncturing of the scrotal sac for access of the vas, and then plug or injected plug occlusion. The success of the aforementioned vas occlusion techniques is not clear and data are still limited. Studies have shown, however, that the time to achieve sterility is longer than the more prominent techniques mentioned in the beginning of this article. The satisfaction rate of patients undergoing IVD techniques has a high rate of satisfaction with regard to the surgery experience itself. Recovery Sexual intercourse can usually be resumed in about a week (depending on recovery); however, pregnancy is still possible as long as the sperm count is above zero. Another method of contraception must be relied upon until a sperm count is performed either two months after the vasectomy or after 10–20 ejaculations have occurred. After a vasectomy, contraceptive precautions must be continued until azoospermia is confirmed. Usually two semen analyses at three and four months are necessary to confirm azoospermia. The British Andrological Society has recommended that a single semen analysis confirming azoospermia after sixteen weeks is sufficient. Post-vasectomy, testicles will continue to produce sperm cells. As before vasectomy, unused sperm are reabsorbed by the body. Conceiving after vasectomy In order to allow the possibility of reproduction via artificial insemination after vasectomy, some opt for cryopreservation of sperm before sterilization. Dr Allan Pacey, senior lecturer in andrology at Sheffield University and secretary of the British Fertility Society, notes that those he sees for a vasectomy reversal which has not worked express wishing they had known they could have stored sperm. Pacey notes, "The problem is you're asking a man to foresee a future where he might not necessarily be with his current partner—and that may be quite hard to do when she's sitting next to you." The cost of cryo-preservation (sperm banking) may also be substantially less than alternative vaso-vasectomy procedures, compared to the costs of in-vitro fertilization (IVF) which usually run from $12,000 to $25,000. Sperm can be aspirated from the testicles or the epididymides, and while there is not enough for successful artificial insemination, there is enough to fertilize an ovum by intracytoplasmic sperm injection. This avoids the problem of antisperm antibodies and may result in a faster pregnancy. IVF may be less costly per cycle than reversal in some health-care systems, but a single IVF cycle is often insufficient for conception. Disadvantages include the need for procedures on the woman, and the standard potential side-effects of IVF for both the mother and the child. Vasectomy reversal Vasectomies are not always reversible. There is a surgical procedure to reverse vasectomies using vasovasostomy (a form of microsurgery first performed by Earl Owen in 1971). Vasovasostomy is effective at achieving pregnancy in a variable percentage of cases, and total out-of-pocket costs in the United States are often upwards of $10,000. The typical success rate of pregnancy following a vasectomy reversal is around 55% if performed within 10 years, and drops to around 25% if performed after 10 years. After reversal, sperm counts and motility are usually much lower than pre-vasectomy levels. There is evidence that those who had a vasectomy may produce more abnormal sperm, which may explain why even a mechanically successful reversal does not always restore fertility. The higher rates of aneuploidy and diploidy in the sperm cells of those who have undergone vasectomy reversal may lead to a higher rate of birth defects. Approximately 2% of men who have undergone vasectomy will undergo a reversal within 10 years of the procedure. A small number of vasectomy reversals are also performed in attempts to relieve post-vasectomy pain syndrome. Prevalence Internationally, vasectomy rates are vastly different. While female sterilisation is the most widely used method worldwide, with 223 million women relying on it, only 28 million women rely on their partner's vasectomy. In the world's 69 least developed countries less than 0.1% of males use vasectomies on average. Of 54 African countries, only ten report measurable vasectomy use and only Eswatini, Botswana, and South Africa exceed 0.1% prevalence. In North America and Europe vasectomy usage is on the order of 10% with some countries reaching 20%. Despite its high efficacy, in the United States, vasectomy is utilized less than half the rate of the alternative female tubal ligation. According to the research, vasectomy in the US is least utilized among black and Latino populations, the groups that have the highest rates of female sterilization. New Zealand, in contrast, has higher levels of vasectomy than tubal ligation; 18% of all males, and 25% of all married males, have had a vasectomy. The age cohort with the highest level of vasectomy was 40–49, where 57% of males had taken it up. Canada, the UK, Bhutan and the Netherlands all have similar levels of uptake. History The first human vasectomies were performed in the late 19th century. The procedure was initially used mainly as a treatment for prostate enlargement and as a eugenic method for sterilizing "degenerates". Vasectomy as a method of voluntary birth control began during the Second World War. The first recorded vasectomy was performed on a dog in 1823. The first human vasectomies were performed to treat benign prostatic hyperplasia, or enlargement of the prostate. Castration had sometimes been used as a treatment for this condition in the 1880s, but, given the serious side effects, doctors sought alternative treatments. The first to suggest vasectomy as an alternative to castration may have been James Ewing Mears (in 1890), or Jean Casimir Félix Guyon. The first human vasectomy is thought to have been performed by Reginald Harrison. By 1900, Harrison reported that he had performed more than 100 vasectomies with no adverse outcomes. In the late 1890s, vasectomy also came to be proposed as a eugenic measure for the sterilization of men considered unfit to reproduce. The first case report of vasectomy in the United States was in 1897, by A. J. Ochsner, a surgeon in Chicago, in a paper titled, "Surgical treatment of habitual criminals". He believed vasectomy to be a simple, effective means for stemming the tide of racial degeneration widely believed to be occurring. In 1902, Harry C. Sharp, the surgeon at the Indiana Reformatory, reported that he had sterilized 42 inmates in an effort to both reduce criminal behavior in those individuals and prevent the birth of future criminals. Eugen Steinach (1861–1944), an Austrian physician, believed that a unilateral vasectomy (severing only one of the two vasa deferentia) in older individuals could restore general vigor and sexual potency, shrink enlarged prostates, and cure various ailments by somehow boosting the hormonal output of the vasectomized testicle. This surgery, which became very popular in the 1920s, was undertaken by many wealthy individuals, including Sigmund Freud and W. B. Yeats. Since these operations lacked rigorous controlled trials, any rejuvenating effect was likely due to the placebo effect, and with the later development of synthetic injectable hormones, this operation fell out of vogue. Vasectomy began to be regarded as a method of consensual birth control during the Second World War. The first vasectomy program on a national scale was launched in 1954 in India. In the 1970s, India enacted a coercive sterilization campaign which resulted in millions of vasectomies. Today, India's sterilization program focuses on coercing poor women. The procedure is seldom performed on dogs, with castration remaining the preferred reproductive control option for canines. It is regularly performed on bulls. Society and culture Availability and legality Vasectomy costs are (or may be) covered in different countries, as a method of both contraception or population control, with some offering it as a part of a national health insurance. The procedure was generally considered illegal in France until 2001, due to provisions in the Napoleonic Code forbidding "self-mutilation". No French law specifically mentioned vasectomy until a 2001 law on contraception and infanticide permitted the procedure. The U.S. Affordable Care Act (signed into law in 2010) does not cover vasectomies, although eight states require state-health insurance plans to cover the cost. These include: Illinois, Maryland, New Jersey, New Mexico, New York, Oregon, Vermont and Washington. In 2014, the Iranian parliament voted for a bill that would ban the procedure. Impact of legal changes (US) An analysis of medical records of 217 million people in the U.S. compared tubal sterilization and vasectomy rates in the last six months of 2021 with rates in the last six months of 2022—just after the Dobbs ruling (i.e. the overturning of Roe vs Wade) in June 2022. Although the effect of Dobbs was different in various social groups, it had a strong impact on those under age 30 with their vasectomy rates increasing by 59%, and tubal sterilization rates increasing by 29%. Tourism Medical tourism, where a patient travels to a less-developed location where a procedure is cheaper to save money and combine convalescence with a vacation, is infrequently used for vasectomy due to its low cost, but is more likely to be used for vasectomy reversal. Many hospitals list vasectomy as being available. Medical tourism has been scrutinized by some governments for quality of care and postoperative care issues. Shooting of Andrew Rynne In 1990, Andrew Rynne, chairperson of the Irish Family Planning Association, and the Republic of Ireland's first vasectomy specialist, was shot by a former client, but he survived. The incident is the subject of a short film, The Vasectomy Doctor, by Paul Webster.	Biology and health sciences	Surgery	Health
69656	https://en.wikipedia.org/wiki/Reactive%20armour	Reactive armour	Reactive armour is a type of vehicle armour used in protecting vehicles, especially modern tanks, against shaped charges and hardened kinetic energy penetrators. The most common type is explosive reactive armour (ERA), but variants include self-limiting explosive reactive armour (SLERA), non-energetic reactive armour (NERA), non-explosive reactive armour (NxRA), and electric armour. NERA and NxRA modules can withstand multiple hits, unlike ERA and SLERA. When a shaped charge strikes the upper plate of the armour, it detonates the inner explosive, releasing blunt damage that the tank can absorb. Reactive armour is intended to counteract anti-tank munitions that work by piercing the armour and then either killing the crew inside, disabling vital mechanical systems, or creating spalling that disables the crew—or all three. Reactive armour can be defeated with multiple hits in the same place, as by tandem-charge weapons, which fire two or more shaped charges in rapid succession. Without tandem charges, hitting precisely the same spot twice is much more difficult. History The Australians were the first recorded to have conceptualized and developed methods to disrupt and spread the jet of a hollow charge shell to reduce its penetrating power. In a June 1944 report from the Explosives Manufacturing Practices Laboratory of the Explosives Factory Maribyrnong, an operational requirement for the defence against shaped charges was laid out. The focus was in regard to Japanese 75 mm hollow charge shells used against Allied tanks in the Pacific. The destructive effect of the shaped charge was identified as caused by a jet moving at high velocities, consisting of particles from the liner. The two methods developed were to destroy the jet by forcing it to act through a layer of explosives, disrupting the jet, and to make it act through a layer of oxidiser, destroying the jet by burning it with oxidising agents. The earliest trials were done with small charges able to defeat 2 inch of steel plate which were readily defeated by a layer of explosive (Baratol, R.D.X., Cordite, etc.) or a vigorous oxidising medium. Subsequent trials with British No.68 and American M9A1 grenades were carried out. However trials were done in few numbers which caused varied results. A mixture of Sodium and Potassium Nitrates explosives was seen as the most practical option due to their casting properties. The mixture acted as an oxidiser which may explode when dispersed and heated. The Explosives Manufacturing Practices Laboratory seemingly developed a more middle road between chemical armor and explosive reactive armor concepts to counter the hollow charge threat. The idea of counterexplosion (kontrvzryv in Russian) in armour was proposed in the USSR by the Scientific Research Institute of Steel (NII Stali) in 1949 by academician Bogdan Vjacheslavovich Voitsekhovsky. The first pre-production models were produced during the 1960s. However, insufficient theoretical analysis during one of the tests resulted in all of the prototype elements being detonated. For a number of reasons, including the aforementioned accident and a belief that Soviet tanks had sufficient armour, the research was ended. No more research was conducted until 1974, when the Ministry of the Defensive Industry announced a contest to find the best tank protection . Picatinny Arsenal, an American military research and manufacturing facility experimented with testing linear cutting charges against anti-tank ammunition in the 1950s, and concluded that they may be effective with an adequate sensing and triggering mechanism, but noted "tactical limitations"; the report was declassified in 1980. A West German researcher, Manfred Held, carried out similar work with the IDF in 1967–1969. Reactive armour created on the basis of the joint research was first installed on Israeli tanks during the 1982 Lebanon war and was judged very effective. Explosive reactive armour An element of explosive reactive armour (ERA) is made of either a sheet or slab of high explosive sandwiched between two metal plates, or multiple "banana shaped" rods filled with high explosive which are referred to as shaped charges. On attack by a penetrating weapon, the explosive detonates, forcibly driving the metal plates apart to damage the penetrator. The shaped charges, in contrast, each detonate individually, launching one spike-shaped plate each, meant to deflect, detonate or cut the incoming projectile. The disruption is attributed to two mechanisms. First, the moving plates change the effective velocity and angle of impact of the shaped charge jet, reducing the angle of incidence and increasing the effective jet velocity versus the plate element. Second, since the plates are angled compared to the usual impact direction of shaped charge warheads, as the plates move outwards the impact point on the plate moves over time, requiring the jet to cut through fresh plates of material. This second effect greatly increases the effective plate thickness during the impact. To be effective against kinetic energy projectiles, ERA must use much thicker and heavier plates and a correspondingly thicker explosive layer. Such heavy ERA, such as the Soviet-developed Kontakt-5, can break apart a penetrating rod that is longer than the ERA is deep, again reducing penetration capability. Such ERA is ineffective against modern APFSDS projectiles, however, due to their depleted uranium construction. An important aspect of ERA is the brisance, or detonation speed of its explosive element. A more brisant explosive and greater plate velocity will result in more plate material being fed into the path of the oncoming jet, greatly increasing the plate's effective thickness. This effect is especially pronounced in the rear plate receding away from the jet, which triples in effective thickness with double the velocity. ERA also counters explosively forged projectiles, as produced by a shaped charge. The counter-explosion must disrupt the incoming projectile so that its momentum is distributed in all directions rather than toward the target, greatly reducing its effectiveness. Explosive reactive armour has been valued by the Soviet Union and its now-independent component states since the 1980s, and almost every tank in the eastern-European military inventory today has either been manufactured to use ERA or had ERA tiles added to it, including even the T-55 and T-62 tanks built forty to fifty years ago, but still used today by reserve units. The U.S. Army uses reactive armour on its Abrams tanks as part of the TUSK (Tank Urban Survivability Kit) package and on Bradley vehicles and the Israelis use it frequently on their American built M60 tanks. ERA tiles are used as add-on (or appliqué) armour to the portions of an armoured fighting vehicle that are most likely to be hit, typically the front (glacis) of the hull and the front and sides of the turret. Their use requires that a vehicle be fairly heavily armoured to protect itself and its crew from the exploding ERA. A further complication to the use of ERA is the inherent danger to anyone near the tank when a plate detonates, though a high-explosive anti-tank (HEAT) warhead explosion would already cause great danger to anyone near the tank. Although ERA plates are intended only to bulge following detonation, the combined energy of the ERA explosive, coupled with the kinetic or explosive energy of the projectile, will frequently cause the plate to explode, creating shrapnel that risks injuring or killing bystanders. Thus, infantry must operate some distance from vehicles protected by ERA in combined arms operations. Sensitivity ERA is insensitive to impact by kinetic projectiles up to 30 mm in caliber. A 20 mm APIT autocannon round penetrates a Serbian ERA sample but fails to detonate it. However, computer simulations indicate that a small caliber (30 mm) HEAT projectile will detonate an ERA, as would larger shape charges and APFSDS penetrators. Non-explosive and non-energetic reactive armour NERA and NxRA operate similarly to explosive reactive armour, but without the explosive liner. Two metal plates sandwich an inert liner, such as rubber. When struck by a shaped charge's metal jet, some of the impact energy is dissipated into the inert liner layer, and the resulting high pressure causes a localized bending or bulging of the plates in the area of the impact. As the plates bulge, the point of jet impact shifts with the plate bulging, increasing the effective thickness of the armour. This is almost the same as the second mechanism that explosive reactive armour uses, but it uses energy from the shaped charge jet rather than from explosives. Since the inner liner is non-explosive, the bulging is less energetic than on explosive reactive armour, and thus offers less protection than a similarly-sized ERA. However, NERA and NxRA are lighter, safe to handle, and safer for nearby infantry; can theoretically be placed on any part of the vehicle; and can be packaged in multiple spaced layers if needed. A key advantage of this kind of armour is that it cannot be defeated by tandem warhead shaped charges, which employ a small forward warhead to detonate ERA before the main warhead fires. Electric armour Electric armour or electromagnetic armour is a proposed reactive armour technology. It is made up of two or more conductive plates separated by an air gap or by an insulating material, creating a high-power capacitor. In operation, a high-voltage power source charges the armour. When an incoming body penetrates the plates, it closes the circuit to discharge the capacitor, dumping energy into the penetrator, which may vaporize it or even turn it into a plasma, diffusing the attack. It is not public knowledge whether this is supposed to function against both kinetic energy penetrators and shaped charge jets, or only the latter. As of 2005, this technology had not yet been introduced on any known operational platform. Another electromagnetic alternative to ERA uses layers of plates of electromagnetic metal with silicone spacers on alternate sides. The damage to the exterior of the armour passes electricity into the plates, causing them to magnetically move together. As the process is completed at the speed of electricity the plates are moving when struck by the projectile, causing the projectile energy to be deflected whilst the energy is also dissipated in parting the magnetically attracted plates.	Technology	Armour	null
69688	https://en.wikipedia.org/wiki/Sterilization%20%28medicine%29	Sterilization (medicine)	Sterilization (also spelled sterilisation) is any of a number of medical methods of permanent birth control that intentionally leaves a person unable to reproduce. Sterilization methods include both surgical and non-surgical options for both males and females. Sterilization procedures are intended to be permanent; reversal is generally difficult. There are multiple ways of having sterilization done, but the two that are used most frequently are tubal ligation for women and vasectomy for men. There are many different ways tubal sterilization can be accomplished. It is extremely effective and in the United States surgical complications are low. With that being said, tubal sterilization is still a method that involves surgery, so there is still a danger. Women that chose a tubal sterilization may have a higher risk of serious side effects, more than a man has with a vasectomy. Pregnancies after a tubal sterilization can still occur, even many years after the procedure. It is not very likely, but if it does happen there is a high risk of ectopic gestation. Statistics confirm that a handful of tubal sterilization surgeries are performed shortly after a vaginal delivery mostly by minilaparotomy. In some cases, sterilization can be reversed but not all. It can vary by the type of sterilization performed. Methods Surgical Surgical sterilization methods include: Tubal ligation in females, known popularly as "having one's tubes tied". The fallopian tubes, which allow the sperm to fertilize the ovum and would carry the fertilized ovum to the uterus, are closed. This generally involves a general anesthetic and a laparotomy or laparoscopic approach to cut, clip or cauterize the fallopian tubes. Bilateral salpingectomy in females, also known as tubal removal. Both fallopian tubes are surgically removed. When done for contraceptive purposes, the ovaries are left in place. This method is considered more effective than tubal ligation, as there is no chance of tubal reconnection or clip failure, and also prevents cancer of the fallopian tubes and can reduce risk of ovarian cancer. Vasoligation in males. The vasa deferentia, the tubes that connect the testicles to the prostate, are cut and closed. This prevents sperm produced in the testicles from entering the ejaculated semen (which is mostly produced in the seminal vesicles and prostate). Although the term vasectomy is established in the general community, the correct medical terminology is vasoligation. Hysterectomy in females. The uterus is surgically removed, permanently preventing pregnancy and some diseases, such as uterine cancer. Castration in males. The testicles are surgically removed. This is frequently used for the sterilization of animals, but rarely for humans. It was also formerly used on some human male children for other reasons; see castrato and eunuch. Transluminal Transluminal procedures are performed by entry through the female reproductive tract. These generally use a catheter to place a substance into the fallopian tubes that eventually causes blockage of the tract in this segment. Such procedures are generally called non-surgical as they use natural orifices and thereby do not necessitate any surgical incision. The Essure procedure was one such transluminal sterilization technique. In this procedure, polyethylene terephthalate fiber inserts were placed into the fallopian tubes, eventually inducing scarring and occlusion of the tubes. In April 2018, the FDA restricted the sale and use of Essure. On July 20, 2018, Bayer announced the halt of sales in the US by the end of 2018. Quinacrine has also been used for transluminal sterilization, but despite a multitude of clinical studies on the use of quinacrine and female sterilization, no randomized, controlled trials have been reported to date and there is some controversy over its use.	Biology and health sciences	Medical procedures	null
69720	https://en.wikipedia.org/wiki/Gastrointestinal%20tract	Gastrointestinal tract	The gastrointestinal tract (GI tract, digestive tract, alimentary canal) is the tract or passageway of the digestive system that leads from the mouth to the anus. The GI tract contains all the major organs of the digestive system, in humans and other animals, including the esophagus, stomach, and intestines. Food taken in through the mouth is digested to extract nutrients and absorb energy, and the waste expelled at the anus as feces. Gastrointestinal is an adjective meaning of or pertaining to the stomach and intestines. Most animals have a "through-gut" or complete digestive tract. Exceptions are more primitive ones: sponges have small pores (ostia) throughout their body for digestion and a larger dorsal pore (osculum) for excretion, comb jellies have both a ventral mouth and dorsal anal pores, while cnidarians and acoels have a single pore for both digestion and excretion. The human gastrointestinal tract consists of the esophagus, stomach, and intestines, and is divided into the upper and lower gastrointestinal tracts. The GI tract includes all structures between the mouth and the anus, forming a continuous passageway that includes the main organs of digestion, namely, the stomach, small intestine, and large intestine. The complete human digestive system is made up of the gastrointestinal tract plus the accessory organs of digestion (the tongue, salivary glands, pancreas, liver and gallbladder). The tract may also be divided into foregut, midgut, and hindgut, reflecting the embryological origin of each segment. The whole human GI tract is about nine meters (30 feet) long at autopsy. It is considerably shorter in the living body because the intestines, which are tubes of smooth muscle tissue, maintain constant muscle tone in a halfway-tense state but can relax in spots to allow for local distention and peristalsis. The gastrointestinal tract contains the gut microbiota, with some 1,000 different strains of bacteria having diverse roles in the maintenance of immune health and metabolism, and many other microorganisms. Cells of the GI tract release hormones to help regulate the digestive process. These digestive hormones, including gastrin, secretin, cholecystokinin, and ghrelin, are mediated through either intracrine or autocrine mechanisms, indicating that the cells releasing these hormones are conserved structures throughout evolution. Human gastrointestinal tract Structure The structure and function can be described both as gross anatomy and as microscopic anatomy or histology. The tract itself is divided into upper and lower tracts, and the intestines small and large parts. Upper gastrointestinal tract The upper gastrointestinal tract consists of the mouth, pharynx, esophagus, stomach, and duodenum. The exact demarcation between the upper and lower tracts is the suspensory muscle of the duodenum. This differentiates the embryonic borders between the foregut and midgut, and is also the division commonly used by clinicians to describe gastrointestinal bleeding as being of either "upper" or "lower" origin. Upon dissection, the duodenum may appear to be a unified organ, but it is divided into four segments based on function, location, and internal anatomy. The four segments of the duodenum are as follows (starting at the stomach, and moving toward the jejunum): bulb, descending, horizontal, and ascending. The suspensory muscle attaches the superior border of the ascending duodenum to the jejunum. The suspensory muscle is an important anatomical landmark that shows the formal division between the duodenum and the jejunum, the first and second parts of the small intestine, respectively. This is a thin muscle which is derived from the embryonic mesoderm. Lower gastrointestinal tract The lower gastrointestinal tract includes most of the small intestine and all of the large intestine. In human anatomy, the intestine (bowel or gut; Greek: éntera) is the segment of the gastrointestinal tract extending from the pyloric sphincter of the stomach to the anus and as in other mammals, consists of two segments: the small intestine and the large intestine. In humans, the small intestine is further subdivided into the duodenum, jejunum, and ileum while the large intestine is subdivided into the cecum, ascending, transverse, descending, and sigmoid colon, rectum, and anal canal. Small intestine The small intestine begins at the duodenum and is a tubular structure, usually between 6 and 7 m long. Its mucosal area in an adult human is about . The combination of the circular folds, the villi, and the microvilli increases the absorptive area of the mucosa about 600-fold, making a total area of about for the entire small intestine. Its main function is to absorb the products of digestion (including carbohydrates, proteins, lipids, and vitamins) into the bloodstream. There are three major divisions: Duodenum: A short structure (about 20–25 cm long) that receives chyme from the stomach, together with pancreatic juice containing digestive enzymes and bile from the gall bladder. The digestive enzymes break down proteins, and bile emulsifies fats into micelles. The duodenum contains Brunner's glands which produce a mucus-rich alkaline secretion containing bicarbonate. These secretions, in combination with bicarbonate from the pancreas, neutralize the stomach acids contained in the chyme. Jejunum: This is the midsection of the small intestine, connecting the duodenum to the ileum. It is about long and contains the circular folds also known as plicae circulares and villi that increase its surface area. Products of digestion (sugars, amino acids, and fatty acids) are absorbed into the bloodstream here. Ileum: The final section of the small intestine. It is about 3 m long, and contains villi similar to the jejunum. It absorbs mainly vitamin B12 and bile acids, as well as any other remaining nutrients. Large intestine The large intestine, also called the colon, forms an arch starting at the cecum and ending at the rectum and anal canal. It also includes the appendix, which is attached to the cecum. Its length is about 1.5 m, and the area of the mucosa in an adult human is about . Its main function is to absorb water and salts. The colon is further divided into: Cecum (first portion of the colon) and appendix Ascending colon (ascending in the back wall of the abdomen) Right colic flexure (flexed portion of the ascending and transverse colon apparent to the liver) Transverse colon (passing below the diaphragm) Left colic flexure (flexed portion of the transverse and descending colon apparent to the spleen) Descending colon (descending down the left side of the abdomen) Sigmoid colon (a loop of the colon closest to the rectum) Rectum Anal canal Development The gut is an endoderm-derived structure. At approximately the sixteenth day of human development, the embryo begins to fold ventrally (with the embryo's ventral surface becoming concave) in two directions: the sides of the embryo fold in on each other and the head and tail fold toward one another. The result is that a piece of the yolk sac, an endoderm-lined structure in contact with the ventral aspect of the embryo, begins to be pinched off to become the primitive gut. The yolk sac remains connected to the gut tube via the vitelline duct. Usually, this structure regresses during development; in cases where it does not, it is known as Meckel's diverticulum. During fetal life, the primitive gut is gradually patterned into three segments: foregut, midgut, and hindgut. Although these terms are often used in reference to segments of the primitive gut, they are also used regularly to describe regions of the definitive gut as well. Each segment of the gut is further specified and gives rise to specific gut and gut-related structures in later development. Components derived from the gut proper, including the stomach and colon, develop as swellings or dilatations in the cells of the primitive gut. In contrast, gut-related derivatives — that is, those structures that derive from the primitive gut but are not part of the gut proper, in general, develop as out-pouchings of the primitive gut. The blood vessels supplying these structures remain constant throughout development. Histology The gastrointestinal tract has a form of general histology with some differences that reflect the specialization in functional anatomy. The GI tract can be divided into four concentric layers in the following order: Mucosa Submucosa Muscular layer Adventitia or serosa Mucosa The mucosa is the innermost layer of the gastrointestinal tract. The mucosa surrounds the lumen, or open space within the tube. This layer comes in direct contact with digested food (chyme). The mucosa is made up of: Epithelium – innermost layer. Responsible for most digestive, absorptive and secretory processes. Lamina propria – a layer of connective tissue. Unusually cellular compared to most connective tissue Muscularis mucosae – a thin layer of smooth muscle that aids the passing of material and enhances the interaction between the epithelial layer and the contents of the lumen by agitation and peristalsis The mucosae are highly specialized in each organ of the gastrointestinal tract to deal with the different conditions. The most variation is seen in the epithelium. Submucosa The submucosa consists of a dense irregular layer of connective tissue with large blood vessels, lymphatics, and nerves branching into the mucosa and muscularis externa. It contains the submucosal plexus, an enteric nervous plexus, situated on the inner surface of the muscularis externa. Muscular layer The muscular layer consists of an inner circular layer and a longitudinal outer layer. The circular layer prevents food from traveling backward and the longitudinal layer shortens the tract. The layers are not truly longitudinal or circular, rather the layers of muscle are helical with different pitches. The inner circular is helical with a steep pitch and the outer longitudinal is helical with a much shallower pitch. Whilst the muscularis externa is similar throughout the entire gastrointestinal tract, an exception is the stomach which has an additional inner oblique muscular layer to aid with grinding and mixing of food. The muscularis externa of the stomach is composed of the inner oblique layer, middle circular layer, and the outer longitudinal layer. Between the circular and longitudinal muscle layers is the myenteric plexus. This controls peristalsis. Activity is initiated by the pacemaker cells, (myenteric interstitial cells of Cajal). The gut has intrinsic peristaltic activity (basal electrical rhythm) due to its self-contained enteric nervous system. The rate can be modulated by the rest of the autonomic nervous system. The coordinated contractions of these layers is called peristalsis and propels the food through the tract. Food in the GI tract is called a bolus (ball of food) from the mouth down to the stomach. After the stomach, the food is partially digested and semi-liquid, and is referred to as chyme. In the large intestine, the remaining semi-solid substance is referred to as faeces. Adventitia and serosa The outermost layer of the gastrointestinal tract consists of several layers of connective tissue. Intraperitoneal parts of the GI tract are covered with serosa. These include most of the stomach, first part of the duodenum, all of the small intestine, caecum and appendix, transverse colon, sigmoid colon and rectum. In these sections of the gut, there is a clear boundary between the gut and the surrounding tissue. These parts of the tract have a mesentery. Retroperitoneal parts are covered with adventitia. They blend into the surrounding tissue and are fixed in position. For example, the retroperitoneal section of the duodenum usually passes through the transpyloric plane. These include the esophagus, pylorus of the stomach, distal duodenum, ascending colon, descending colon and anal canal. In addition, the oral cavity has adventitia. Gene and protein expression Approximately 20,000 protein coding genes are expressed in human cells and 75% of these genes are expressed in at least one of the different parts of the digestive organ system. Over 600 of these genes are more specifically expressed in one or more parts of the GI tract and the corresponding proteins have functions related to digestion of food and uptake of nutrients. Examples of specific proteins with such functions are pepsinogen PGC and the lipase LIPF, expressed in chief cells, and gastric ATPase ATP4A and gastric intrinsic factor GIF, expressed in parietal cells of the stomach mucosa. Specific proteins expressed in the stomach and duodenum involved in defence include mucin proteins, such as mucin 6 and intelectin-1. Transit time The time taken for food to transit through the gastrointestinal tract varies on multiple factors, including age, ethnicity, and gender. Several techniques have been used to measure transit time, including radiography following a barium-labeled meal, breath hydrogen analysis, scintigraphic analysis following a radiolabeled meal, and simple ingestion and spotting of corn kernels. It takes 2.5 to 3 hours for 50% of the contents to leave the stomach. The rate of digestion is also dependent of the material being digested, as food composition from the same meal may leave the stomach at different rates. Total emptying of the stomach takes around 4–5 hours, and transit through the colon takes 30 to 50 hours. Immune function The gastrointestinal tract forms an important part of the immune system. Immune barrier The surface area of the digestive tract is estimated to be about 32 square meters, or about half a badminton court. With such a large exposure (more than three times larger than the exposed surface of the skin), these immune components function to prevent pathogens from entering the blood and lymph circulatory systems. Fundamental components of this protection are provided by the intestinal mucosal barrier, which is composed of physical, biochemical, and immune elements elaborated by the intestinal mucosa. Microorganisms also are kept at bay by an extensive immune system comprising the gut-associated lymphoid tissue (GALT) There are additional factors contributing to protection from pathogen invasion. For example, low pH (ranging from 1 to 4) of the stomach is fatal for many microorganisms that enter it. Similarly, mucus (containing IgA antibodies) neutralizes many pathogenic microorganisms. Other factors in the GI tract contribution to immune function include enzymes secreted in the saliva and bile. Immune system homeostasis Beneficial bacteria also can contribute to the homeostasis of the gastrointestinal immune system. For example, Clostridia, one of the most predominant bacterial groups in the GI tract, play an important role in influencing the dynamics of the gut's immune system. It has been demonstrated that the intake of a high fiber diet could be responsible for the induction of T-regulatory cells (Tregs). This is due to the production of short-chain fatty acids during the fermentation of plant-derived nutrients such as butyrate and propionate. Basically, the butyrate induces the differentiation of Treg cells by enhancing histone H3 acetylation in the promoter and conserved non-coding sequence regions of the FOXP3 locus, thus regulating the T cells, resulting in the reduction of the inflammatory response and allergies. Intestinal microbiota The large intestine contains multiple types of bacteria that can break down molecules the human body cannot process alone, demonstrating a symbiotic relationship. These bacteria are responsible for gas production at host–pathogen interface, which is released as flatulence. Intestinal bacteria can also participate in biosynthesis reactions. For example, certain strains in the large intestine produce vitamin B12; an essential compound in humans for things like DNA synthesis and red blood cell production. However, the primary function of the large intestine is water absorption from digested material (regulated by the hypothalamus) and the reabsorption of sodium and nutrients. Beneficial intestinal bacteria compete with potentially harmful bacteria for space and "food", as the intestinal tract has limited resources. A ratio of 80–85% beneficial to 15–20% potentially harmful bacteria is proposed for maintaining homeostasis. An imbalanced ratio results in dysbiosis. Detoxification and drug metabolism Enzymes such as CYP3A4, along with the antiporter activities, are also instrumental in the intestine's role of drug metabolism in the detoxification of antigens and xenobiotics. Other animals In most vertebrates, including amphibians, birds, reptiles, egg-laying mammals, and some fish, the gastrointestinal tract ends in a cloaca and not an anus. In the cloaca, the urinary system is fused with the genito-anal pore. Therians (all mammals that do not lay eggs, including humans) possess separate anal and uro-genital openings. The females of the subgroup Placentalia have even separate urinary and genital openings. During early development, the asymmetric position of the bowels and inner organs is initiated (see also axial twist theory). Ruminants show many specializations for digesting and fermenting tough plant material, consisting of additional stomach compartments. Many birds and other animals have a specialised stomach in the digestive tract called a gizzard used for grinding up food. Another feature found in a range of animals is the crop. In birds this is found as a pouch alongside the esophagus. In 2020, the oldest known fossil digestive tract, of an extinct wormlike organism in the Cloudinidae was discovered; it lived during the late Ediacaran period about 550 million years ago. A through-gut (one with both mouth and anus) is thought to have evolved within the nephrozoan clade of Bilateria, after their ancestral ventral orifice (single, as in cnidarians and acoels; re-evolved in nephrozoans like flatworms) stretched antero-posteriorly, before the middle part of the stretch would get narrower and closed fully, leaving an anterior orifice (mouth) and a posterior orifice (anus plus genital opening). A stretched gut without the middle part closed is present in another branch of bilaterians, the extinct proarticulates. This and the amphistomic development (when both mouth and anus develop from the gut stretch in the embryo) present in some nephrozoans (e.g. roundworms) are considered to support this hypothesis. Clinical significance Diseases There are many diseases and conditions that can affect the gastrointestinal system, including infections, inflammation and cancer. Various pathogens, such as bacteria that cause foodborne illnesses, can induce gastroenteritis which results from inflammation of the stomach and small intestine. Antibiotics to treat such bacterial infections can decrease the microbiome diversity of the gastrointestinal tract, and further enable inflammatory mediators. Gastroenteritis is the most common disease of the GI tract. Gastrointestinal cancer may occur at any point in the gastrointestinal tract, and includes mouth cancer, tongue cancer, oesophageal cancer, stomach cancer, and colorectal cancer. Inflammatory conditions. Ileitis is an inflammation of the ileum, colitis is an inflammation of the large intestine. Appendicitis is inflammation of the appendix located at the caecum. This is a potentially fatal condition if left untreated; most cases of appendicitis require surgical intervention. Diverticular disease is a condition that is very common in older people in industrialized countries. It usually affects the large intestine but has been known to affect the small intestine as well. Diverticulosis occurs when pouches form on the intestinal wall. Once the pouches become inflamed it is known as diverticulitis. Inflammatory bowel disease is an inflammatory condition affecting the bowel walls, and includes the subtypes Crohn's disease and ulcerative colitis. While Crohn's can affect the entire gastrointestinal tract, ulcerative colitis is limited to the large intestine. Crohn's disease is widely regarded as an autoimmune disease. Although ulcerative colitis is often treated as though it were an autoimmune disease, there is no consensus that it actually is such. Functional gastrointestinal disorders the most common of which is irritable bowel syndrome. Functional constipation and chronic functional abdominal pain are other functional disorders of the intestine that have physiological causes but do not have identifiable structural, chemical, or infectious pathologies. Symptoms Several symptoms can indicate problems with the gastrointestinal tract, including: Vomiting, which may include regurgitation of food or the vomiting of blood Diarrhea, or the passage of liquid or more frequent stools Constipation, which refers to the passage of fewer and hardened stools Blood in stool, which includes fresh red blood, maroon-coloured blood, and tarry-coloured blood Treatment Gastrointestinal surgery can often be performed in the outpatient setting. In the United States in 2012, operations on the digestive system accounted for 3 of the 25 most common ambulatory surgery procedures and constituted 9.1 percent of all outpatient ambulatory surgeries. Imaging Various methods of imaging the gastrointestinal tract include the upper and lower gastrointestinal series: Radioopaque dyes may be swallowed to produce a barium swallow Parts of the tract may be visualised by camera. This is known as endoscopy if examining the upper gastrointestinal tract and colonoscopy or sigmoidoscopy if examining the lower gastrointestinal tract. Capsule endoscopy is where a capsule containing a camera is swallowed in order to examine the tract. Biopsies may also be taken when examined. An abdominal X-ray may be used to examine the lower gastrointestinal tract. Other related diseases Cholera Enteric duplication cyst Giardiasis Pancreatitis Peptic ulcer disease Yellow fever Helicobacter pylori is a gram-negative spiral bacterium. Over half the world's population is infected with it, mainly during childhood; it is not certain how the disease is transmitted. It colonizes the gastrointestinal system, predominantly the stomach. The bacterium has specific survival conditions that are specific to the human gastric microenvironment: it is both capnophilic and microaerophilic. Helicobacter also exhibits a tropism for gastric epithelial lining and the gastric mucosal layer about it. Gastric colonization of this bacterium triggers a robust immune response leading to moderate to severe inflammation, known as gastritis. Signs and symptoms of infection are gastritis, burning abdominal pain, weight loss, loss of appetite, bloating, burping, nausea, bloody vomit, and black tarry stools. Infection can be detected in a number of ways: GI X-rays, endoscopy, blood tests for anti-Helicobacter antibodies, a stool test, and a urease breath test (which is a by-product of the bacteria). If caught soon enough, it can be treated with three doses of different proton pump inhibitors as well as two antibiotics, taking about a week to cure. If not caught soon enough, surgery may be required. Intestinal pseudo-obstruction is a syndrome caused by a malformation of the digestive system, characterized by a severe impairment in the ability of the intestines to push and assimilate. Symptoms include daily abdominal and stomach pain, nausea, severe distension, vomiting, heartburn, dysphagia, diarrhea, constipation, dehydration and malnutrition. There is no cure for intestinal pseudo-obstruction. Different types of surgery and treatment managing life-threatening complications such as ileus and volvulus, intestinal stasis which lead to bacterial overgrowth, and resection of affected or dead parts of the gut may be needed. Many patients require parenteral nutrition. Ileus is a blockage of the intestines. Coeliac disease is a common form of malabsorption, affecting up to 1% of people of northern European descent. An autoimmune response is triggered in intestinal cells by digestion of gluten proteins. Ingestion of proteins found in wheat, barley and rye, causes villous atrophy in the small intestine. Lifelong dietary avoidance of these foodstuffs in a gluten-free diet is the only treatment. Enteroviruses are named by their transmission-route through the intestine (enteric meaning intestinal), but their symptoms are not mainly associated with the intestine. Endometriosis can affect the intestines, with similar symptoms to IBS. Bowel twist (or similarly, bowel strangulation) is a comparatively rare event (usually developing sometime after major bowel surgery). It is, however, hard to diagnose correctly, and if left uncorrected can lead to bowel infarction and death. (The singer Maurice Gibb is understood to have died from this.) Angiodysplasia of the colon Constipation Diarrhea Hirschsprung's disease (aganglionosis) Intussusception Polyp (medicine) (see also colorectal polyp) Pseudomembranous colitis Toxic megacolon usually a complication of ulcerative colitis Uses of animal guts Intestines from animals other than humans are used in a number of ways. From each species of livestock that is a source of milk, a corresponding rennet is obtained from the intestines of milk-fed . Pig and calf intestines are eaten, and pig intestines are used as sausage casings. Calf intestines supply calf-intestinal alkaline phosphatase (CIP), and are used to make goldbeater's skin. Other uses are: The use of animal gut strings by musicians can be traced back to the third dynasty of Egypt. In the recent past, strings were made out of lamb gut. With the advent of the modern era, musicians have tended to use strings made of silk, or synthetic materials such as nylon or steel. Some instrumentalists, however, still use gut strings in order to evoke the older tone quality. Although such strings were commonly referred to as "catgut" strings, cats were never used as a source for gut strings. Sheep gut was the original source for natural gut string used in racquets, such as for tennis. Today, synthetic strings are much more common, but the best gut strings are now made out of cow gut. Gut cord has also been used to produce strings for the snares that provide a snare drum's characteristic buzzing timbre. While the modern snare drum almost always uses metal wire rather than gut cord, the North African bendir frame drum still uses gut for this purpose. "Natural" sausage hulls, or casings, are made of animal gut, especially hog, beef, and lamb. The wrapping of kokoretsi, gardoubakia, and torcinello is made of lamb (or goat) gut. Haggis is traditionally boiled in, and served in, a sheep stomach. Chitterlings, a kind of food, consist of thoroughly washed pig's gut. Animal gut was used to make the cord lines in longcase clocks and for fusee movements in bracket clocks, but may be replaced by metal wire. The oldest known condoms, from 1640 AD, were made from animal intestine.	Biology and health sciences	Human anatomy	null
69785	https://en.wikipedia.org/wiki/Pelican	Pelican	Pelicans (genus Pelecanus) are a genus of large water birds that make up the family Pelecanidae. They are characterized by a long beak and a large throat pouch used for catching prey and draining water from the scooped-up contents before swallowing. They have predominantly pale plumage, except for the brown and Peruvian pelicans. The bills, pouches, and bare facial skin of all pelicans become brightly coloured before the breeding season. The eight living pelican species have a patchy, seasonally-dependent yet global distribution, ranging latitudinally from the tropics to the temperate zone. Pelicans are absent from interior Amazonian South America, from polar regions and the open ocean; at least one species is known to migrate to the inland desert of Australia's Red Centre, after heavy rains create temporary lakes. White pelicans are also observed at the American state of Utah's Great Salt Lake, for example, some 600 miles (965 km) from the nearest coastline (the Pacific West Coast). They have also been seen hundreds of miles inland in North America, having flown northwards along the Mississippi River and other large waterways. Long thought to be related to frigatebirds, cormorants, tropicbirds, and gannets and boobies, pelicans instead are most closely related to the shoebill and hamerkop storks (although these two birds are not actually true 'storks'), and are placed in the order Pelecaniformes. Ibises, spoonbills, herons, and bitterns have been classified in the same order. Fossil evidence of pelicans dates back at least 36 million years to the remains of a tibiotarsus recovered from late Eocene strata of Egypt that bears striking similarity to modern species of pelican. They are thought to have evolved in the Old World and spread into the Americas; this is reflected in the relationships within the genus as the eight species divide into Old World and New World lineages. This hypothesis is supported by fossil evidence from the oldest pelican taxa. Pelicans will frequent inland waterways but are most known for residing along maritime and coastal zones, where they feed principally on fish in their large throat pouches, diving into the water and catching them at/near the water's surface. They can adapt to varying degrees of water salinity, from freshwater and brackish to—most commonly—seawater. They are gregarious birds, travelling in flocks, hunting cooperatively, and breeding colonially. Four white-plumaged species tend to nest on the ground, and four brown or grey-plumaged species nest mainly in trees. The relationship between pelicans and people has often been contentious. The birds have been persecuted because of their perceived competition with commercial and recreational fishing. Their populations have fallen through habitat destruction, disturbance, and environmental pollution, and three species are of conservation concern. They also have a long history of cultural significance in mythology, and in Christian and heraldic iconography. Taxonomy and systematics Etymology The name comes from the Ancient Greek word pelekan (πελεκάν), which is itself derived from the word pelekys (πέλεκυς) meaning "axe". In classical times, the word was applied to both the pelican and the woodpecker. The genus Pelecanus was first formally described by Carl Linnaeus in his landmark 1758 10th edition of Systema Naturae. He described the distinguishing characteristics as a straight bill hooked at the tip, linear nostrils, a bare face, and fully webbed feet. This early definition included frigatebirds, cormorants, and sulids, as well as pelicans. Taxonomy The family Pelecanidae was introduced (as Pelicanea) by the French polymath Constantine Samuel Rafinesque in 1815. Pelicans give their name to the Pelecaniformes, an order which has a varied taxonomic history. Tropicbirds, darters, cormorants, gannets, boobies, and frigatebirds, all traditional members of the order, have since been reclassified: tropicbirds into their own order, Phaethontiformes, and the remainder into the Suliformes. In their place, herons, ibises, spoonbills, the hamerkop, and the shoebill have now been transferred into the Pelecaniformes. Molecular evidence suggests that the shoebill and the hamerkop form a sister group to the pelicans, though some doubt exists as to the exact relationships among the three lineages. The oldest known record of Pelicans is a right tibiotarsus very similar to those of modern species from the Birket Qarun Formation in the Wadi El Hitan in Egypt, dating to the late Eocene (Priabonian), referred to the genus Eopelecanus. Living species The eight living pelican species were traditionally divided into two groups, one containing four ground-nesters with mainly white adult plumage (Australian, Dalmatian, great white, and American white pelicans), and one containing four grey- or brown-plumaged species which nest preferentially either in trees (pink-backed, spot-billed and brown pelicans), or on sea rocks (Peruvian pelican). The largely marine brown and Peruvian pelicans, formerly considered conspecific, are sometimes separated from the others by placement in the subgenus Leptopelecanus but in fact species with both sorts of appearance and nesting behavior are found in either. DNA sequencing of both mitochondrial and nuclear genes yielded quite different relationships; the three New World pelicans formed one lineage, with the American white pelican sister to the two brown pelicans, and the five Old World species the other. The Dalmatian, pink-backed, and spot-billed were all closely related to one another, while the Australian white pelican was their next-closest relative. The great white pelican also belonged to this lineage, but was the first to diverge from the common ancestor of the other four species. This finding suggests that pelicans evolved in the Old World and spread into the Americas, and that preference for tree- or ground-nesting is more related to size than genetics. Fossil record The fossil record shows that the pelican lineage has existed for at least 36 million years; the oldest known pelican fossil was assigned to Eopelecanus aegyptiacus and was found in late Eocene (middle to late part of the early Priabonian stage/age) deposits of the Birket Qarun Formation within the Wadi Al-Hitan World Heritage Site in Egypt. A more complete fossil pelican of early Oligocene age is known from deposits at the Luberon in southeastern France, and is remarkably similar to modern forms. Its beak is almost complete and is morphologically identical to that of present-day pelicans, showing that this advanced feeding apparatus was already in existence at the time. An Early Miocene fossil has been named Miopelecanus gracilis on the basis of certain features originally considered unique, but later thought to lie within the range of interspecific variation in Pelecanus. The Late Eocene Protopelicanus may be a pelecaniform or suliform – or a similar aquatic bird such as a pseudotooth (Pelagornithidae), but is not generally considered a pelecanid. The supposed Miocene pelican Liptornis from Patagonia is a nomen dubium (of doubtful validity), being based on fragments providing insufficient evidence to support a valid description. Fossil finds from North America have been meagre compared with Europe, which has a richer fossil record. Several Pelecanus species have been described from fossil material, including: Pelecanus cadimurka, Rich & van Tets, 1981 (Late Pliocene, South Australia) Pelecanus cautleyi, Davies, 1880 (Early Pliocene, Siwalik Hills, India) Pelecanus fraasi, Lydekker, 1891 (Middle Miocene, Bavaria, Germany) Pelecanus gracilis, Milne-Edwards, 1863 (Early Miocene, France) (see: Miopelecanus) Pelecanus halieus, Wetmore, 1933 (Late Pliocene, Idaho, US) Pelecanus intermedius, Fraas, 1870 (Middle Miocene, Bavaria, Germany) (transferred to Miopelecanus by Cheneval in 1984) Pelecanus odessanus, Widhalm, 1886 (Late Miocene, near Odesa, Ukraine) Pelecanus paranensis, Noriega et al., 2023 (Late Miocene, Entre Ríos Province, Argentina) Pelecanus schreiberi, Olson, 1999 (Early Pliocene, North Carolina, US) Pelecanus sivalensis, Davies, 1880 (Early Pliocene, Siwalik Hills, India) Pelecanus tirarensis, Miller, 1966 (Late Oligocene to Middle Miocene, South Australia) Description Pelicans are very large birds with very long bills characterised by a downcurved hook at the end of the upper mandible, and the attachment of a huge gular pouch to the lower. The slender rami of the lower bill and the flexible tongue muscles form the pouch into a basket for catching fish, and sometimes rainwater, though to not hinder the swallowing of large fish, the tongue itself is tiny. They have a long neck and short stout legs with large, fully webbed feet. Although they are among the heaviest of flying birds, they are relatively light for their apparent bulk because of air pockets in the skeleton and beneath the skin, enabling them to float high in the water. The tail is short and square. The wings are long and broad, suitably shaped for soaring and gliding flight, and have the unusually large number of 30 to 35 secondary flight feathers. Males are generally larger than females and have longer bills. The smallest species is the brown pelican, small individuals of which can be no more than and long, with a wingspan of as little as . The largest is believed to be the Dalmatian, at up to and in length, with a maximum wingspan of . The Australian pelican's bill may grow up to long in large males, the longest of any bird. Pelicans have mainly light-coloured plumage, the exceptions being the brown and Peruvian pelicans. The bills, pouches, and bare facial skin of all species become brighter before breeding season commences. The throat pouch of the Californian subspecies of the brown pelican turns bright red, and fades to yellow after the eggs are laid, while the throat pouch of the Peruvian pelican turns blue. The American white pelican grows a prominent knob on its bill that is shed once females have laid eggs. The plumage of immature pelicans is darker than that of adults. Newly hatched chicks are naked and pink, darkening to grey or black after four to 14 days, then developing a covering of white or grey down. Air sacs Anatomical dissections of two brown pelicans in 1939 showed that pelicans have a network of air sacs under their skin situated across the ventral surface including the throat, breast, and undersides of the wings, as well as having air sacs in their bones. The air sacs are connected to the airways of the respiratory system, and the pelican can keep its air sacs inflated by closing its glottis, but how air sacs are inflated is not clear. The air sacs serve to keep the pelican remarkably buoyant in the water and may also cushion the impact of the pelican's body on the water surface when they dive from flight into water to catch fish. Superficial air sacs may also help to round body contours (especially over the abdomen, where surface protuberances may be caused by viscera changing size and position) to enable the overlying feathers to form more effective heat insulation and also to enable feathers to be held in position for good aerodynamics. Distribution and habitat Modern pelicans are found on all continents except Antarctica. They primarily inhabit warm regions, although breeding ranges extend to latitudes of 45° South (Australian pelicans in Tasmania) and 60° North (American white pelicans in western Canada). Birds of inland and coastal waters, they are absent from polar regions, the deep ocean, oceanic islands (except the Galapagos), and inland South America, as well as from the eastern coast of South America from the mouth of the Amazon River southwards. Subfossil bones have been recovered from as far south as New Zealand's South Island, although their scarcity and isolated occurrence suggests that these remains may have merely been vagrants from Australia (much as is the case today). Behaviour and ecology Pelicans swim well with their strong legs and their webbed feet. They rub the backs of their heads on their preen glands to pick up an oily secretion, which they transfer to their plumage to waterproof it. Holding their wings only loosely against their bodies, pelicans float with relatively little of their bodies below the water surface. They dissipate excess heat by gular flutter – rippling the skin of the throat and pouch with the bill open to promote evaporative cooling. They roost and loaf communally on beaches, sandbanks, and in shallow water. A fibrous layer deep in the breast muscles can hold the wings rigidly horizontal for gliding and soaring. Thus, they use thermals for soaring to heights of 3,000 m (10,000 ft) or more, combined both with gliding and with flapping flight in V formation, to commute distances up to to feeding areas. Pelicans also fly low (or "skim") over stretches of water, using a phenomenon known as ground effect to reduce drag and increase lift. As the air flows between the wings and the water surface, it is compressed to a higher density and exerts a stronger upward force against the bird above. Hence, substantial energy is saved while flying. Adult pelicans rely on visual displays and behaviour to communicate, particularly using their wings and bills. Agonistic behaviour consists of thrusting and snapping at opponents with their bills, or lifting and waving their wings in a threatening manner. Adult pelicans grunt when at the colony, but are generally silent elsewhere or outside breeding season. Conversely, colonies are noisy, as chicks vocalise extensively. Breeding and lifespan Pelicans are gregarious and nest colonially. Pairs are monogamous for a single season, but the pair bond extends only to the nesting area; mates are independent away from the nest. The ground-nesting (white) species have a complex communal courtship involving a group of males chasing a single female in the air, on land, or in the water while pointing, gaping, and thrusting their bills at each other. They can finish the process in a day. The tree-nesting species have a simpler process in which perched males advertise for females. The location of the breeding colony is constrained by the availability of an ample supply of fish to eat, although pelicans can use thermals to soar and commute for hundreds of kilometres daily to fetch food. The Australian pelican has two reproductive strategies depending on the local degree of environmental predictability. Colonies of tens or hundreds, rarely thousands, of birds breed regularly on small coastal and subcoastal islands where food is seasonally or permanently available. In arid inland Australia, especially in the endorheic Lake Eyre basin, pelicans breed opportunistically in very large numbers of up to 50,000 pairs, when irregular major floods, which may be many years apart, fill ephemeral salt lakes and provide large amounts of food for several months before drying out again. In all species, copulation takes place at the nest site; it begins shortly after pairing and continues for three to ten days before egg-laying. The male brings the nesting material, in ground-nesting species (which may not build a nest) sometimes in the pouch, and in tree-nesting species crosswise in the bill. The female then heaps the material up to form a simple structure. The eggs are oval, white, and coarsely textured. All species normally lay at least two eggs; the usual clutch size is one to three, rarely up to six. Both sexes incubate with the eggs on top of or below the feet; they may display when changing shifts. Incubation takes 30–36 days; hatching success for undisturbed pairs can be as high as 95%, but because of sibling competition or siblicide, in the wild, usually all but one nestling dies within the first few weeks (later in the pink-backed and spot-billed species). Both parents feed their young. Small chicks are fed by regurgitation; after about a week, they are able to put their heads into their parents' pouches and feed themselves. Sometimes before, but especially after being fed the pelican chick may seem to "throw a tantrum" by loudly vocalizing and dragging itself around in a circle by one wing and leg, striking its head on the ground or anything nearby and the tantrums sometimes end in what looks like a seizure that results in the chick falling briefly unconscious; the reason is not clearly known, but a common belief is that it is to draw attention to itself and away from any siblings who are waiting to be fed. Parents of ground-nesting species sometimes drag older young around roughly by the head before feeding them. From about 25 days old, the young of these species gather in "pods" or "crèches" of up to 100 birds in which parents recognise and feed only their own offspring. By six to eight weeks, they wander around, occasionally swimming, and may practise communal feeding. Young of all species fledge ten to 12 weeks after hatching. They may remain with their parents afterwards, but are now seldom or never fed. They are mature at three or four years old. Overall breeding success is highly variable. Pelicans live for 15 to 25 years in the wild, although one reached an age of 54 years in captivity. Feeding The diet of pelicans usually consists of fish, but occasionally amphibians, turtles, crustaceans, insects, birds, and mammals are also eaten. The size of the preferred prey fish varies depending on pelican species and location. For example, in Africa, the pink-backed pelican generally takes fish ranging in size from fry up to and the great white pelican prefers somewhat larger fish, up to , but in Europe, the latter species has been recorded taking fish up to . In deep water, white pelicans often fish alone. Nearer the shore, several encircle schools of small fish or form a line to drive them into the shallows, beating their wings on the water surface and then scooping up the prey. Although all pelican species may feed in groups or alone, the Dalmatian, pink-backed, and spot-billed pelicans are the only ones to prefer solitary feeding. When fishing in groups, all pelican species have been known to work together to catch their prey, and Dalmatian pelicans may even cooperate with great cormorants. Large fish are caught with the bill-tip, then tossed up in the air to be caught and slid into the gullet head-first. A gull will sometimes stand on the pelican's head, peck it to distraction, and grab a fish from the open bill. Pelicans in their turn sometimes snatch prey from other waterbirds. The brown pelican usually plunge-dives head-first for its prey, from a height as great as , especially for anchovies and menhaden. The only other pelican to feed using a similar technique is the Peruvian pelican, but its dives are typically from a lower height than the brown pelican. The Australian and American white pelicans may feed by low plunge-dives landing feet-first and then scooping up the prey with the beak, but they—as well as the remaining pelican species—primarily feed while swimming on the water. Aquatic prey is most commonly taken at or near the water surface. Although principally a fish eater, the Australian pelican is also an eclectic and opportunistic scavenger and carnivore that forages in landfill sites, as well as taking carrion and "anything from insects and small crustaceans to ducks and small dogs". Food is not stored in a pelican's throat pouch, contrary to popular folklore. Pelicans may also eat birds. In southern Africa, eggs and chicks of the Cape cormorant are an important food source for great white pelicans. Several other bird species have been recorded in the diet of this pelican in South Africa, including Cape gannet chicks on Malgas Island as well as crowned cormorants, kelp gulls, greater crested terns, and African penguins on Dassen Island and elsewhere. The Australian pelican, which is particularly willing to take a wide range of prey items, has been recorded feeding on young Australian white ibis, and young and adult grey teals and silver gulls. Brown pelicans have been reported preying on young common murres in California and the eggs and nestlings of cattle egrets and nestling great egrets in Baja California, Mexico. Peruvian pelicans in Chile have been recorded feeding on nestlings of imperial shags, juvenile Peruvian diving petrels, and grey gulls. Cannibalism of chicks of their own species is known from the Australian, brown, and Peruvian pelicans. Non-native great white pelicans have been observed swallowing city pigeons in St. James's Park in London, England. Status and conservation Populations Globally, pelican populations are adversely affected by these main factors: declining supplies of fish through overfishing or water pollution, destruction of habitat, direct effects of human activity such as disturbance at nesting colonies, hunting and culling, entanglement in fishing lines and hooks, and the presence of pollutants such as DDT and endrin. Most species' populations are more or less stable, although three are classified by the IUCN as being at risk. All species breed readily in zoos, which is potentially useful for conservation management. The combined population of brown and Peruvian pelicans is estimated at 650,000 birds, with around 250,000 in the United States and Caribbean, and 400,000 in Peru. The National Audubon Society estimates the global population of the brown pelican at 300,000. Numbers of brown pelican plummeted in the 1950s and 1960s, largely as a consequence of environmental DDT pollution, and the species was listed as endangered in the US in 1970. With restrictions on DDT use in the US from 1972, its population has recovered, and it was delisted in 2009. The Peruvian pelican is listed as near threatened because, although the population is estimated by BirdLife International to exceed 500,000 mature individuals, and is possibly increasing, it has been much higher in the past. It declined dramatically during the 1998 El Niño event and could experience similar declines in the future. Conservation needs include regular monitoring throughout the range to determine population trends, particularly after El Niño years, restricting human access to important breeding colonies, and assessing interactions with fisheries. The spot-billed pelican has an estimated population between 13,000 and 18,000 and is considered to be near threatened in the IUCN Red List of Threatened Species. Numbers declined substantially during the 20th century, one crucial factor being the eradication of the important Sittaung valley breeding colony in Burma through deforestation and the loss of feeding sites. The chief threats it faces are from habitat loss and human disturbance, but populations have mostly stabilised following increased protection in India and Cambodia. The pink-backed pelican has a large population ranging over much of sub-Saharan Africa. In the absence of substantial threats or evidence of declines across its range, its conservation status is assessed as being of least concern. Regional threats include the drainage of wetlands and increasing disturbance in southern Africa. The species is susceptible to bioaccumulation of toxins and the destruction of nesting trees by logging. The American white pelican has increased in numbers, with its population estimated at over 157,000 birds in 2005, becoming more numerous east of the continental divide, while declining in the west. However, whether its numbers have been affected by exposure to pesticides is unclear, as it has also lost habitat through wetland drainage and competition with recreational use of lakes and rivers. Great white pelicans range over a large area of Africa and southern Asia. The overall trend in numbers is uncertain, with a mix of regional populations that are increasing, declining, stable, or unknown, but no evidence has been found of rapid overall decline, and the status of the species is assessed as being of least concern. Threats include the drainage of wetlands, persecution and sport hunting, disturbance at the breeding colonies, and contamination by pesticides and heavy metals. The Dalmatian pelican has a population estimated at between 10,000 and 20,000 following massive declines in the 19th and 20th centuries. The main ongoing threats include hunting, especially in eastern Asia, disturbance, coastal development, collision with overhead power lines, and the over-exploitation of fish stocks. It is listed as near threatened by the IUCN Red List of Threatened Species as the population trend is downwards, especially in Mongolia, where it is nearly extinct. However, several European colonies are increasing in size and the largest colony for the species, at the Small Prespa Lake in Greece, has reached about 1,400 breeding pairs following conservation measures. Widespread across Australia, the Australian pelican has a population generally estimated at between 300,000 and 500,000 individuals. Overall population numbers fluctuate widely and erratically depending on wetland conditions and breeding success across the continent. The species is assessed as being of least concern. Culling and disturbance Pelicans have been persecuted by humans for their perceived competition for fish, despite the fact that their diet overlaps little with fish caught by people. Starting in the 1880s, American white pelicans were clubbed and shot, their eggs and young were deliberately destroyed, and their feeding and nesting sites were degraded by water management schemes and wetland drainage. Even in the 21st century, an increase in the population of American white pelicans in southeastern Idaho in the US was seen to threaten the recreational cutthroat trout fishery there, leading to official attempts to reduce pelican numbers through systematic harassment and culling. Great white pelicans on Dyer Island, in the Western Cape region of South Africa, were culled during the 19th century because their predation of the eggs and chicks of guano-producing seabirds was seen to threaten the livelihood of the guano collectors. More recently, such predation at South African seabird colonies has impacted on the conservation of threatened seabird populations, especially crowned cormorants, Cape cormorants, and bank cormorants. This has led to suggestions that pelican numbers should be controlled at vulnerable colonies. Apart from habitat destruction and deliberate, targeted persecution, pelicans are vulnerable to disturbance at their breeding colonies by birdwatchers, photographers, and other curious visitors. Human presence alone can cause the birds to accidentally displace or destroy their eggs, leave hatchlings exposed to predators and adverse weather, or even abandon their colonies completely. Poisoning and pollution DDT pollution in the environment was a major cause of decline of brown pelican populations in North America in the 1950s and 1960s. It entered the oceanic food web, contaminating and accumulating in several species, including one of the pelican's primary food fish – the northern anchovy. Its metabolite DDE is a reproductive toxicant in pelicans and many other birds, causing eggshell thinning and weakening, and consequent breeding failure through the eggs being accidentally crushed by brooding birds. Since an effective ban on the use of DDT was implemented in the US in 1972, the eggshells of breeding brown pelicans there have thickened and their populations have largely recovered. In the late 1960s, following the major decline in brown pelican numbers in Louisiana from DDT poisoning, 500 pelicans were imported from Florida to augment and re-establish the population; over 300 subsequently died in April and May 1975 from poisoning by the pesticide endrin. About 14,000 pelicans, including 7,500 American white pelicans, perished from botulism after eating fish from the Salton Sea in 1990. In 1991, abnormal numbers of brown pelicans and Brandt's cormorants died at Santa Cruz, California, when their food fish (anchovies) were contaminated with neurotoxic domoic acid, produced by the diatom Pseudo-nitzschia. As waterbirds that feed on fish, pelicans are highly susceptible to oil spills, both directly by being oiled and by the impact on their food resources. A 2007 report to the California Fish and Game Commission estimated that during the previous 20 years, some 500–1,000 brown pelicans had been affected by oil spills in California. A 2011 report by the Center for Biological Diversity, a year after the April 2010 Deepwater Horizon oil spill, said that 932 brown pelicans had been collected after being affected by oiling and estimated that ten times that number had been harmed as a result of the spill. Where pelicans interact with fishers, through either sharing the same waters or scavenging for fishing refuse, they are especially vulnerable to being hooked and entangled in both active and discarded fishing lines. Fish hooks are swallowed or catch in the skin of the pouch or webbed feet, and strong monofilament fishing line can become wound around bill, wings, or legs, resulting in crippling, starvation, and often death. Local rescue organisations have been established in North America and Australia by volunteers to treat and rehabilitate injured pelicans and other wildlife. Parasites and disease As with other bird families, pelicans are susceptible to a variety of parasites. Avian malaria is carried by the mosquito Culex pipens, and high densities of these biting insects may force pelican colonies to be abandoned. Leeches may attach to the vent or sometimes the inside of the pouch. A study of the parasites of the American white pelican found 75 different species, including tapeworms, flukes, flies, fleas, ticks, and nematodes. The brown pelican has a similarly extensive range of parasites. The nematodes Contracaecum multipapillatum and C. mexicanum and the trematode Ribeiroia ondatrae have caused illness and mortality in the Puerto Rican population, possibly endangering the pelican on this island. Many pelican parasites are found in other bird groups, but several lice are very host-specific. Healthy pelicans can usually cope with their lice, but sick birds may carry hundreds of individuals, which hastens a sick bird's demise. The pouch louse Piagetiella peralis occurs in the pouch and so it cannot be removed by preening. While this is usually not a serious problem even when present in such numbers that it covers the whole interior of the pouch, sometimes inflammation and bleeding may occur from it and harm the host. In May 2012, hundreds of Peruvian pelicans were reported to have perished in Peru from a combination of starvation and roundworm infestation. Religion, mythology, and popular culture The pelican (henet in Egyptian) was associated in Ancient Egypt with death and the afterlife. It was depicted in art on the walls of tombs, and figured in funerary texts, as a protective symbol against snakes. Henet was also referred to in the Pyramid Texts as the "mother of the king" and thus seen as a goddess.	Biology and health sciences	Pelecanimorphae	null
69794	https://en.wikipedia.org/wiki/Cryptogram	Cryptogram	A cryptogram is a type of puzzle that consists of a short piece of encrypted text. Generally the cipher used to encrypt the text is simple enough that the cryptogram can be solved by hand. Substitution ciphers where each letter is replaced by a different letter, number, or symbol are frequently used. To solve the puzzle, one must recover the original lettering. Though once used in more serious applications, they are now mainly printed for entertainment in newspapers and magazines. Other types of classical ciphers are sometimes used to create cryptograms. An example is the book cipher, where a book or article is used to encrypt a message. History The ciphers used in cryptograms were created not for entertainment purposes, but for real encryption of military or personal secrets. The first use of the cryptogram for entertainment purposes occurred during the Middle Ages by monks who had spare time for intellectual games. A manuscript found at Bamberg states that Irish visitors to the court of Merfyn Frych ap Gwriad (died 844), king of Gwynedd in Wales, were given a cryptogram which could only be solved by transposing the letters from Latin into Greek. Around the thirteenth century, the English monk Roger Bacon wrote a book in which he listed seven cipher methods, and stated that "a man is crazy who writes a secret in any other way than one which will conceal it from the vulgar." In the 19th century Edgar Allan Poe helped to popularize cryptograms with many newspaper and magazine articles. Well-known examples of cryptograms in contemporary culture are the syndicated newspaper puzzles Cryptoquip and Cryptoquote, from King Features. Celebrity Cipher, distributed by Andrew McMeel, is another cipher game in contemporary culture, challenging the player to decrypt quotes from famous personalities. A cryptoquip is a specific type of cryptogram that usually comes with a clue or a pun. The solution often involves a humorous or witty phrase. In a public challenge, writer J.M. Appel announced on September 28, 2014, that the table of contents page of his short story collection, Scouting for the Reaper, doubled as a cryptogram, and he pledged an award for the first to solve it. Solving a cryptogram Cryptograms based on substitution ciphers can often be solved by frequency analysis and by recognizing letter patterns in words, such as one-letter words, which, in English, can only be "i" or "a" (and sometimes "o"). Double letters, apostrophes, and the fact that no letter can substitute for itself in the cipher also offer clues to the solution. Occasionally, cryptogram puzzle makers will start the solver off with a few letters. A printed code key form; the alphabet with a blank under each letter to fill in the substituted letter, is usually not provided but can be drawn to use as a solving aid if needed. Skilled puzzle solvers should require neither a code key form nor starter clue letters. Other crypto puzzles While the cryptogram has remained popular, over time other puzzles similar to it have emerged. One of these is the Cryptoquote, which is a famous quote encrypted in the same way as a cryptogram. A more recent version, with a biblical twist, is CodedWord. This puzzle makes the solution available only online, where it provides a short exegesis on the biblical text. A third is the Cryptoquiz. The top of this puzzle has a category (unencrypted), such as "Flowers". Below this is a list of encrypted words which are related to the stated category. The person must then solve for the entire list to finish the puzzle. Yet another type involves using numbers as they relate to texting to solve the puzzle. The Zodiac Killer sent four cryptograms to police while he was still active. Despite much research, only two of these have been translated, which was of no help in identifying the serial killer.	Technology	Computer security	null
69893	https://en.wikipedia.org/wiki/Headache	Headache	A headache, also known as cephalalgia, is the symptom of pain in the face, head, or neck. It can occur as a migraine, tension-type headache, or cluster headache. There is an increased risk of depression in those with severe headaches. Headaches can occur as a result of many conditions. There are a number of different classification systems for headaches. The most well-recognized is that of the International Headache Society, which classifies it into more than 150 types of primary and secondary headaches. Causes of headaches may include dehydration; fatigue; sleep deprivation; stress; the effects of medications (overuse) and recreational drugs, including withdrawal; viral infections; loud noises; head injury; rapid ingestion of a very cold food or beverage; and dental or sinus issues (such as sinusitis). Treatment of a headache depends on the underlying cause, but commonly involves pain medication (especially in case of migraine or cluster headaches). A headache is one of the most commonly experienced of all physical discomforts. About half of adults have a headache in a given year. Tension headaches are the most common, affecting about 1.6 billion people (21.8% of the population) followed by migraine headaches which affect about 848 million (11.7%). Causes There are more than 200 types of headaches. Some are harmless and some are life-threatening. The description of the headache and findings on neurological examination, determine whether additional tests are needed and what treatment is best. Headaches are broadly classified as "primary" or "secondary". Primary headaches are benign, recurrent headaches not caused by underlying disease or structural problems. For example, migraine is a type of primary headache. While primary headaches may cause significant daily pain and disability, they are not dangerous from a physiological point of view. Secondary headaches are caused by an underlying disease, like an infection, head injury, vascular disorders, brain bleed, stomach irritation, or tumors. Secondary headaches can be dangerous. Certain "red flags" or warning signs indicate a secondary headache may be dangerous. Primary Ninety percent of all headaches are primary headaches. Primary headaches usually first start when people are between 20 and 40 years old. The most common types of primary headaches are migraines and tension-type headaches. They have different characteristics. Migraines typically present with pulsing head pain, nausea, photophobia (sensitivity to light) and phonophobia (sensitivity to sound). Tension-type headaches usually present with non-pulsing "bandlike" pressure on both sides of the head, not accompanied by other symptoms. Such kind of headaches may be further classified into-episodic and chronic tension type headaches Other very rare types of primary headaches include: cluster headaches: short episodes (15–180 minutes) of severe pain, usually around one eye, with autonomic symptoms (tearing, red eye, nasal congestion) which occur at the same time every day. Cluster headaches can be treated with triptans and prevented with prednisone, ergotamine or lithium. trigeminal neuralgia or occipital neuralgia: shooting face pain hemicrania continua: continuous unilateral pain with episodes of severe pain. Hemicrania continua can be relieved by the medication indomethacin. primary stabbing headache: recurrent episodes of stabbing "ice pick pain" or "jabs and jolts" for 1 second to several minutes without autonomic symptoms (tearing, red eye, nasal congestion). These headaches can be treated with indomethacin. primary cough headache: starts suddenly and lasts for several minutes after coughing, sneezing or straining (anything that may increase pressure in the head). Serious causes (see secondary headaches red flag section) must be ruled out before a diagnosis of "benign" primary cough headache can be made. primary exertional headache: throbbing, pulsatile pain which starts during or after exercising, lasting for 5 minutes to 24 hours. The mechanism behind these headaches is unclear, possibly due to straining causing veins in the head to dilate, causing pain. These headaches can be prevented by not exercising too strenuously and can be treated with medications such as indomethacin. primary sex headache: dull, bilateral headache that starts during sexual activity and becomes much worse during orgasm. These headaches are thought to be due to lower pressure in the head during sex. It is important to realize that headaches that begin during orgasm may be due to a subarachnoid hemorrhage, so serious causes must be ruled out first. These headaches are treated by advising the person to stop sex if they develop a headache. Medications such as propranolol and diltiazem can also be helpful. hypnic headache: a moderate-severe headache that starts a few hours after falling asleep and lasts 15–30 minutes. The headache may recur several times during the night. Hypnic headaches are usually in older women. They may be treated with lithium. Secondary Headaches may be caused by problems elsewhere in the head or neck. Some of these are not harmful, such as cervicogenic headache (pain arising from the neck muscles). The excessive use of painkillers can paradoxically cause worsening painkiller headaches. More serious causes of secondary headaches include the following: meningitis: inflammation of the meninges which presents with fever and meningismus, or stiff neck ischemic stroke or a previous stage of the same hemorragic stroke or a previous stage of the same intracranial hemorrhage (bleeding inside the brain) because of any origin subarachnoid hemorrhage (with acute, severe headache, stiff neck without fever) because of any origin intraparenchymal hemorrhage (with headache only) because of any origin ruptured aneurysm or aneurysm brain tumor (a form of cancer): dull headache, worse with exertion and change in position, accompanied by nausea and vomiting. Often, the person will have nausea and vomiting for weeks before the headache starts. temporal arteritis: inflammatory disease of arteries common in the elderly (average age 70) with fever, headache, weight loss, jaw claudication, tender vessels by the temples, polymyalgia rheumatica acute closed-angle glaucoma (increased pressure in the eyeball): a headache that starts with eye pain, blurry vision, associated with nausea and vomiting. On physical exam, the person will have red eyes and a fixed, mid-dilated pupil. arteriovenous malformation post-ictal headaches: Headaches that happen after a convulsion or other type of seizure, as part of the period after the seizure (the post-ictal state) Gastrointestinal disorders may cause headaches, including Helicobacter pylori infection, celiac disease, non-celiac gluten sensitivity, irritable bowel syndrome, inflammatory bowel disease, gastroparesis, and hepatobiliary disorders. The treatment of the gastrointestinal disorders may lead to a remission or improvement of headaches. Migraine headaches are also associated with Cyclic Vomiting Syndrome (CVS). CVS is characterized by episodes of severe vomiting, and often occur alongside symptoms similar to those of migraine headaches (photophobia, abdominal pain, etc.). Pathophysiology The brain itself is not sensitive to pain, because it lacks pain receptors. However, several areas of the head and neck do have pain receptors and can thus sense pain. These include the extracranial arteries, middle meningeal artery, large veins, venous sinuses, cranial and spinal nerves, head and neck muscles, the meninges, falx cerebri, parts of the brainstem, eyes, ears, teeth, and lining of the mouth. Pial arteries, rather than pial veins are responsible for pain production. Headaches often result from traction or irritation of the meninges and blood vessels. The pain receptors may be stimulated by head trauma or tumours and cause headaches. Blood vessel spasms, dilated blood vessels, inflammation or infection of meninges and muscular tension can also stimulate pain receptors. Once stimulated, a nociceptor sends a message up the length of the nerve fibre to the nerve cells in the brain, signalling that a part of the body hurts. Primary headaches are more difficult to understand than secondary headaches. The exact mechanisms which cause migraines, tension headaches and cluster headaches are not known. There have been different hypotheses over time that attempt to explain what happens in the brain to cause these headaches. Migraines are currently thought to be caused by dysfunction of the nerves in the brain. Previously, migraines were thought to be caused by a primary problem with the blood vessels in the brain. This vascular theory, which was developed in the 20th century by Wolff, suggested that the aura in migraines is caused by constriction of intracranial vessels (vessels inside the brain), and the headache itself is caused by rebound dilation of extracranial vessels (vessels just outside the brain). Dilation of these extracranial blood vessels activates the pain receptors in the surrounding nerves, causing a headache. The vascular theory is no longer accepted. Studies have shown migraine head pain is not accompanied by extracranial vasodilation, but rather only has some mild intracranial vasodilation. Currently, most specialists think migraines are due to a primary problem with the nerves in the brain. Auras are thought to be caused by a wave of increased activity of neurons in the cerebral cortex (a part of the brain) known as cortical spreading depression followed by a period of depressed activity. Some people think headaches are caused by the activation of sensory nerves which release peptides or serotonin, causing inflammation in arteries, dura and meninges and also cause some vasodilation. Triptans, medications that treat migraines, block serotonin receptors and constrict blood vessels. People who are more susceptible to experiencing migraines without headaches are those who have a family history of migraines, women, and women who are experiencing hormonal changes or are taking birth control pills or are prescribed hormone replacement therapy. Tension headaches are thought to be caused by the activation of peripheral nerves in the head and neck muscles. Cluster headaches involve overactivation of the trigeminal nerve and hypothalamus in the brain, but the exact cause is unknown. Diagnosis Most headaches can be diagnosed by the clinical history alone. If the symptoms described by the person sound dangerous, further testing with neuroimaging or lumbar puncture may be necessary. Electroencephalography (EEG) is not useful for headache diagnosis. The first step to diagnosing a headache is to determine if the headache is old or new. A "new headache" can be a headache that has started recently, or a chronic headache that has changed character. For example, if a person has chronic weekly headaches with pressure on both sides of his head, and then develops a sudden severe throbbing headache on one side of his head, they have a new headache. Red flags It can be challenging to differentiate between low-risk, benign headaches and high-risk, dangerous headaches since symptoms are often similar. Headaches that are possibly dangerous require further lab tests and imaging to diagnose. The American College for Emergency Physicians published criteria for low-risk headaches. They are as follows: age younger than 30 years features typical of primary headache history of similar headache no abnormal findings on neurologic exam no concerning change in normal headache pattern no high-risk comorbid conditions (for example, HIV) no new concerning history or physical examination findings A number of characteristics make it more likely that the headache is due to potentially dangerous secondary causes which may be life-threatening or cause long-term damage. These "red flag" symptoms mean that a headache warrants further investigation with neuroimaging and lab tests. In general, people complaining of their "first" or "worst" headache warrant imaging and further workup. People with progressively worsening headache also warrant imaging, as they may have a mass or a bleed that is gradually growing, pressing on surrounding structures and causing worsening pain. People with neurological findings on exam, such as weakness, also need further workup. The American Headache Society recommends using "SSNOOP", a mnemonic to remember the red flags for identifying a secondary headache: Systemic symptoms (fever or weight loss) Systemic disease (HIV infection, malignancy) Neurologic symptoms or signs Onset sudden (thunderclap headache) Onset after age 40 years Previous headache history (first, worst, or different headache) Other red flag symptoms include: Old headaches Old headaches are usually primary headaches and are not dangerous. They are most often caused by migraines or tension headaches. Migraines are often unilateral, pulsing headaches accompanied by nausea or vomiting. There may be an aura (visual symptoms, numbness or tingling) 30–60 minutes before the headache, warning the person of a headache. Migraines may also not have auras. Tension-type headaches usually have bilateral "bandlike" pressure on both sides of the head usually without nausea or vomiting. However, some symptoms from both headache groups may overlap. It is important to distinguish between the two because the treatments are different. The mnemonic 'POUND' helps distinguish between migraines and tension-type headaches. POUND stands for: One review article found that if 4–5 of the POUND characteristics are present, a migraine is 24 times as likely a diagnosis than a tension-type headache (likelihood ratio 24). If 3 characteristics of POUND are present, migraine is 3 times more likely a diagnosis than tension type headache (likelihood ratio 3). If only 2 POUND characteristics are present, tension-type headaches are 60% more likely (likelihood ratio 0.41). Another study found the following factors independently each increase the chance of migraine over tension-type headache: nausea, photophobia, phonophobia, exacerbation by physical activity, unilateral, throbbing quality, chocolate as a headache trigger, and cheese as a headache trigger. Cluster headaches are relatively rare (1 in 1000 people) and are more common in men than women. They present with sudden onset explosive pain around one eye and are accompanied by autonomic symptoms (tearing, runny nose and red eye). Temporomandibular jaw pain (chronic pain in the jaw joint), and cervicogenic headache (headache caused by pain in muscles of the neck) are also possible diagnoses. For chronic, unexplained headaches, keeping a headache diary can be useful for tracking symptoms and identifying triggers, such as association with menstrual cycle, exercise and food. While mobile electronic diaries for smartphones are becoming increasingly common, a recent review found most are developed with a lack of evidence base and scientific expertise. Cephalalgiaphobia is fear of headaches or getting a headache. New headaches New headaches are more likely to be dangerous secondary headaches. They can, however, simply be the first presentation of a chronic headache syndrome, like migraine or a tension headache. One recommended diagnostic approach is as follows. If any urgent red flags are present such as visual impairment, new seizures, new weakness, or new confusion, further workup with imaging and possibly a lumbar puncture should be done (see red flags section for more details). If the headache is sudden onset (thunderclap headache), a computed tomography scan (CT scan) to look for a brain bleed (subarachnoid hemorrhage) should be done. If the CT scan does not show a bleed, a lumbar puncture should be done to look for blood in the cerebrospinal fluid (CSF), as the CT scan can be falsely negative and subarachnoid hemorrhages can be fatal. If there are signs of infection such as fever, rash, or stiff neck, a lumbar puncture to look for meningitis should be considered. In an older person, if there is jaw claudication and scalp tenderness, a temporal artery biopsy should be preformed to look for temporal arteritis, immediate treatment should be started, if results of the biopsy are positive. Neuroimaging Old headaches The US Headache Consortium has guidelines for neuroimaging of non-acute headaches. Most old, chronic headaches do not require neuroimaging. If a person has the characteristic symptoms of a migraine, neuroimaging is not needed as it is very unlikely the person has an intracranial abnormality. If the person has neurological findings, such as weakness, on exam, neuroimaging may be considered. New headaches All people who present with red flags indicating a dangerous secondary headache should receive neuroimaging. The best form of neuroimaging for these headaches is controversial. Non-contrast computerized tomography (CT) scan is usually the first step in head imaging as it is readily available in Emergency Departments and hospitals and is cheaper than MRI. Non-contrast CT is best for identifying an acute head bleed. Magnetic Resonance Imaging (MRI) is best for brain tumors and problems in the posterior fossa, or back of the brain. MRI is more sensitive for identifying intracranial problems, however it can pick up brain abnormalities that are not relevant to the person's headaches. The American College of Radiology recommends the following imaging tests for different specific situations: Lumbar puncture A lumbar puncture is a procedure in which cerebral spinal fluid is removed from the spine with a needle. A lumbar puncture is necessary to look for infection or blood in the spinal fluid. A lumbar puncture can also evaluate the pressure in the spinal column, which can be useful for people with idiopathic intracranial hypertension (usually young, obese women who have increased intracranial pressure), or other causes of increased intracranial pressure. In most cases, a CT scan should be done first. Classification Headaches are most thoroughly classified by the International Headache Society's International Classification of Headache Disorders (ICHD), which published the second edition in 2004. The third edition of the International Headache Classification was published in 2013 in a beta version ahead of the final version. This classification is accepted by the WHO. Other classification systems exist. One of the first published attempts was in 1951. The US National Institutes of Health developed a classification system in 1962. ICHD-2 The International Classification of Headache Disorders (ICHD) is an in-depth hierarchical classification of headaches published by the International Headache Society. It contains explicit (operational) diagnostic criteria for headache disorders. The first version of the classification, ICHD-1, was published in 1988. The current revision, ICHD-2, was published in 2004. The classification uses numeric codes. The top, one-digit diagnostic level includes 14 headache groups. The first four of these are classified as primary headaches, groups 5-12 as secondary headaches, cranial neuralgia, central and primary facial pain and other headaches for the last two groups. The ICHD-2 classification defines migraines, tension-types headaches, cluster headache and other trigeminal autonomic headache as the main types of primary headaches. Also, according to the same classification, stabbing headaches and headaches due to cough, exertion and sexual activity (sexual headache) are classified as primary headaches. The daily-persistent headaches along with the hypnic headache and thunderclap headaches are considered primary headaches as well. Secondary headaches are classified based on their cause and not on their symptoms. According to the ICHD-2 classification, the main types of secondary headaches include those that are due to head or neck trauma such as whiplash injury, intracranial hematoma, post craniotomy or other head or neck injury. Headaches caused by cranial or cervical vascular disorders such as ischemic stroke and transient ischemic attack, non-traumatic intracranial hemorrhage, vascular malformations or arteritis are also defined as secondary headaches. This type of headache may also be caused by cerebral venous thrombosis or different intracranial vascular disorders. Other secondary headaches are those due to intracranial disorders that are not vascular such as low or high pressure of the cerebrospinal fluid pressure, non-infectious inflammatory disease, intracranial neoplasm, epileptic seizure or other types of disorders or diseases that are intracranial but that are not associated with the vasculature of the central nervous system. ICHD-2 classifies headaches that are caused by the ingestion of a certain substance or by its withdrawal as secondary headaches as well. This type of headache may result from the overuse of some medications or exposure to some substances. HIV/AIDS, intracranial infections and systemic infections may also cause secondary headaches. The ICHD-2 system of classification includes the headaches associated with homeostasis disorders in the category of secondary headaches. This means that headaches caused by dialysis, high blood pressure, hypothyroidism, cephalalgia and even fasting are considered secondary headaches. Secondary headaches, according to the same classification system, can also be due to the injury of any of the facial structures including teeth, jaws, or temporomandibular joint. Headaches caused by psychiatric disorders such as somatization or psychotic disorders are also classified as secondary headaches. The ICHD-2 classification puts cranial neuralgias and other types of neuralgia in a different category. According to this system, there are 19 types of neuralgias and headaches due to different central causes of facial pain. Moreover, the ICHD-2 includes a category that contains all the headaches that cannot be classified. Although the ICHD-2 is the most complete headache classification there is and it includes frequency in the diagnostic criteria of some types of headaches (primarily primary headaches), it does not specifically code frequency or severity which are left at the discretion of the examiner. NIH The NIH classification consists of brief definitions of a limited number of headaches. The NIH system of classification is more succinct and only describes five categories of headaches. In this case, primary headaches are those that do not show organic or structural causes. According to this classification, primary headaches can only be vascular, myogenic, cervicogenic, traction, and inflammatory. Management Primary headache syndromes have many different possible treatments. In those with chronic headaches the long term use of opioids appears to result in greater harm than benefit. Secondary headaches (caused by another disease) Treatment of secondary headaches involves treating their underlying cause. For example, a person with meningitis will require antibiotics, and a person with a brain tumor may require surgery, chemotherapy or brain radiation. The possible origins of a headache have been studied and classified. Migraines Migraine can be somewhat improved by lifestyle changes, but most people require medicines to control their symptoms. Medications are either to prevent getting migraines, or to reduce symptoms once a migraine starts. Preventive medications are generally recommended when people have more than four attacks of migraine per month, headaches last longer than 12 hours or the headaches are very disabling. Possible therapies include beta blockers, antidepressants, anticonvulsants and NSAIDs. The type of preventive medicine is usually chosen based on the other symptoms the person has. For example, if the person also has depression, an antidepressant is a good choice. Abortive therapies for migraines may be oral, if the migraine is mild to moderate, or may require stronger medicine given intravenously or intramuscularly. Mild to moderate headaches should first be treated with acetaminophen (paracetamol) or NSAIDs, like ibuprofen. If accompanied by nausea or vomiting, an antiemetic such as metoclopramide (Reglan) can be given orally or rectally. Moderate to severe attacks should be treated first with an oral triptan, a medication that mimics serotonin (an agonist) and causes mild vasoconstriction. If accompanied by nausea and vomiting, parenteral (through a needle in the skin) triptans and antiemetics can be given. Sphenopalatine ganglion block (SPG block, also known nasal ganglion block or pterygopalatine ganglion blocks) can abort and prevent migraines, tension headaches and cluster headaches. It was originally described by American ENT surgeon Greenfield Sluder in 1908. Both blocks and neurostimulation have been studied as treatment for headaches. Several complementary and alternative strategies can help with migraines. The American Academy of Neurology guidelines for migraine treatment in 2000 stated relaxation training, electromyographic feedback and cognitive behavioral therapy may be considered for migraine treatment, along with medications. Tension-type headaches Tension-type headaches can usually be managed with NSAIDs (ibuprofen, naproxen, aspirin), or acetaminophen. Triptans are not helpful in tension-type headaches unless the person also has migraines. For chronic tension type headaches, amitriptyline is the only medication proven to help. Amitriptyline is a medication which treats depression and also independently treats pain. It works by blocking the reuptake of serotonin and norepinephrine, and also reduces muscle tenderness by a separate mechanism. Studies evaluating acupuncture for tension-type headaches have been mixed. Overall, they show that acupuncture is probably not helpful for tension-type headaches. Cluster headaches Abortive therapy for cluster headaches includes subcutaneous sumatriptan (injected under the skin) and triptan nasal sprays. High flow oxygen therapy also helps with relief. For people with extended periods of cluster headaches, preventive therapy can be necessary. Verapamil is recommended as first line treatment. Lithium can also be useful. For people with shorter bouts, a short course of prednisone (10 days) can be helpful. Ergotamine is useful if given 1–2 hours before an attack. Neuromodulation Peripheral neuromodulation has tentative benefits in primary headaches including cluster headaches and chronic migraine. How it may work is still being looked into. Epidemiology Literature reviews find that approximately 64–77% of adults have had a headache at some point in their lives. During each year, on average, 46–53% of people have headaches. However, the prevalence of headache varies widely depending on how the survey was conducted, with studies finding lifetime prevalence of as low as 8% to as high as 96%. Most of these headaches are not dangerous. Only approximately 1–5% of people who seek emergency treatment for headaches have a serious underlying cause. More than 90% of headaches are primary headaches. Most of these primary headaches are tension headaches. Most people with tension headaches have "episodic" tension headaches that come and go. Only 3.3% of adults have chronic tension headaches, with headaches for more than 15 days in a month. Approximately 12–18% of people in the world have migraines. More women than men experience migraines. In Europe and North America, 5–9% of men experience migraines, while 12–25% of women experience migraines. Cluster headaches are relatively uncommon. They affect only 1–3 per thousand people in the world. Cluster headaches affect approximately three times as many men as women. History The first recorded classification system was published by Aretaeus of Cappadocia, a medical scholar of Greco-Roman antiquity. He made a distinction between three different types of headache: i) cephalalgia, by which he indicates a sudden onset, temporary headache; ii) cephalea, referring to a chronic type of headache; and iii) heterocrania, a paroxysmal headache on one side of the head. Another classification system that resembles the modern ones was published by Thomas Willis, in De Cephalalgia in 1672. In 1787 Christian Baur generally divided headaches into idiopathic (primary headaches) and symptomatic (secondary ones), and defined 84 categories. Children In general, children experience the same types of headaches as adults do, but their symptoms may be slightly different. The diagnostic approach to headaches in children is similar to that of adults. However, young children may not be able to verbalize pain well. If a young child is fussy, they may have a headache. Approximately 1% of emergency department visits for children are for headache. Most of these headaches are not dangerous. The most common type of headache seen in pediatric emergency rooms is headache caused by a cold (28.5%). Other headaches diagnosed in the emergency department include post-traumatic headache (20%), headache related to a problem with a ventriculoperitoneal shunt (a device put into the brain to remove excess CSF and reduce pressure in the brain) (11.5%) and migraine (8.5%). The most common serious headaches found in children include brain bleeds (subdural hematoma, epidural hematoma), brain abscesses, meningitis and ventriculoperitoneal shunt malfunction. Only 4–6.9% of kids with a headache have a serious cause. Just as in adults, most headaches are benign, but when head pain is accompanied with other symptoms such as speech problems, muscle weakness, and loss of vision, a more serious underlying cause may exist: hydrocephalus, meningitis, encephalitis, abscess, hemorrhage, tumor, blood clots, or head trauma. In these cases, the headache evaluation may include CT scan or MRI in order to look for possible structural disorders of the central nervous system. If a child with a recurrent headache has a normal physical exam, neuroimaging is not recommended. Guidelines state children with abnormal neurologic exams, confusion, seizures and recent onset of worst headache of life, change in headache type or anything suggesting neurologic problems should receive neuroimaging. When children complain of headaches, many parents are concerned about a brain tumor. Generally, headaches caused by brain masses are incapacitating and accompanied by vomiting. One study found characteristics associated with brain tumor in children are: headache for greater than 6 months, headache related to sleep, vomiting, confusion, no visual symptoms, no family history of migraine and abnormal neurologic exam. Some measures can help prevent headaches in children. Drinking plenty of water throughout the day, avoiding caffeine, getting enough and regular sleep, eating balanced meals at the proper times, and reducing stress and excess of activities may prevent headaches. Treatments for children are similar to those for adults, however certain medications such as narcotics should not be given to children. Children who have headaches will not necessarily have headaches as adults. In one study of 100 children with headache, eight years later 44% of those with tension headache and 28% of those with migraines were headache free. In another study of people with chronic daily headache, 75% did not have chronic daily headaches two years later, and 88% did not have chronic daily headaches eight years later. Cardiac Cephalgia in Heart Attack Cardiac cephalgia is a rare type of headache occurring during myocardial infarction, characterized by sudden, severe head pain that typically develops during or immediately following a heart attack. The pain is usually located in the occipital or frontal regions and can be accompanied by other cardiac symptoms like chest pain, shortness of breath, or radiating arm pain. This specific headache type is considered a potential warning sign of cardiac distress and requires immediate medical attention to prevent potentially life-threatening complications.	Biology and health sciences	Symptoms and signs	Health
69902	https://en.wikipedia.org/wiki/Extreme%20weather	Extreme weather	Extreme weather includes unexpected, unusual, severe, or unseasonal weather; weather at the extremes of the historical distribution—the range that has been seen in the past. Extreme events are based on a location's recorded weather history. They are defined as lying in the most unusual ten percent (10th or 90th percentile of a probability density function). The main types of extreme weather include heat waves, cold waves and heavy precipitation or storm events, such as tropical cyclones. The effects of extreme weather events are economic costs, loss of human lives, droughts, floods, landslides. Severe weather is a particular type of extreme weather which poses risks to life and property. Weather patterns can experience some variation, and so extreme weather can be attributed, at least in part, to the natural climate variability that exists on Earth. For example, the El Niño-Southern Oscillation (ENSO) or the North Atlantic Oscillation (NAO) are climate phenomena that impact weather patterns worldwide. Generally speaking, one event in extreme weather cannot be attributed to any one single cause. However, certain system wide changes to global weather systems can lead to increased frequency or intensity of extreme weather events. Climate change is making some extreme weather events more frequent and more intense. This applies in particular to heat waves and cold waves. The science of extreme event attribution looks at the reasons behind extreme events. Scientists are fairly sure that climate change makes heavy rainfall events as well as drought periods more severe. Climate models indicate that rising temperatures will make extreme weather events worse worldwide. Extreme weather has serious impacts on human society and on ecosystems. There is loss of human lives, damage to infrastructure and ecosystem destruction. For example, a global insurer Munich Re estimates that natural disasters cause more than 90 billion dollars in global direct losses in 2015. Some human activities can exacerbate the effects, for example poor urban planning, wetland destruction, and building homes along floodplains. Definition Extreme weather describes unusual weather events that are at the extremes of the historical distribution for a given area. The IPCC Sixth Assessment Report defines an extreme weather event as follows: "An event that is rare at a particular place and time of year. Definitions of 'rare' vary, but an extreme weather event would normally be as rare as or rarer than the 10th or 90th percentile of a probability density function estimated from observations." In comparison, the term severe weather is any aspect of the weather that poses risks to life, property or requires the intervention of authorities. Severe weather is thus a particular type of extreme weather. Types Definitions of extreme weather vary in different parts of the community, changing the outcomes of research from those fields. Heat waves Heat waves are periods of abnormally high temperatures and heat index. Definitions of a heatwave vary because of the variation of temperatures in different geographic locations. Excessive heat is often accompanied by high levels of humidity, but can also be catastrophically dry. Because heat waves are not visible as other forms of severe weather, like hurricanes, tornadoes, and thunderstorms, they are one of the less known forms of extreme weather. Severely hot weather can damage populations and crops due to potential dehydration or hyperthermia, heat cramps, heat expansion, and heat stroke. Dried soils are more susceptible to erosion, decreasing lands available for agriculture. Outbreaks of wildfires can increase in frequency as dry vegetation has an increased likelihood of igniting. The evaporation of bodies of water can be devastating to marine populations, decreasing the size of the habitats available as well as the amount of nutrition present within the waters. Livestock and other animal populations may decline as well. During excessive heat, plants shut their leaf pores (stomata), a protective mechanism to conserve water but also curtails plants' absorption capabilities. This leaves more pollution and ozone in the air, which leads to higher mortality in the population. It has been estimated that extra pollution during the hot summer of 2006 in the UK, cost 460 lives. The European heat waves from summer 2003 are estimated to have caused 30,000 excess deaths, due to heat stress and air pollution. Over 200 U.S cities have registered new record high temperatures. The worst heat wave in the US occurred in 1936 and killed more than 5000 people directly. The worst heat wave in Australia occurred in 1938–39 and killed 438. The second worst was in 1896. Power outages can also occur within areas experiencing heat waves due to the increased demand for electricity (i.e. air conditioning use). The urban heat island effect can increase temperatures, particularly overnight. Cold waves A cold wave is a weather phenomenon that is distinguished by a cooling of the air. Specifically, as used by the U.S. National Weather Service, a cold wave is a rapid fall in temperature within a 24-hour period requiring substantially increased protection for agriculture, industry, commerce, and social activities. The precise criterion for a cold wave is determined by the rate at which the temperature falls, and the minimum to which it falls. This minimum temperature is dependent on the geographical region and time of year. Cold waves generally are capable of occurring at any geological location and are formed by large cool air masses that accumulate over certain regions, caused by movements of air streams. A cold wave can cause death and injury to livestock and wildlife. Exposure to cold mandates greater caloric intake for all animals, including humans, and if a cold wave is accompanied by heavy and persistent snow, grazing animals may be unable to reach necessary food and water, and die of hypothermia or starvation. Cold waves often necessitate the purchase of fodder for livestock at a considerable cost to farmers. Human populations can be inflicted with frostbite when exposed for extended periods of time to cold and may result in the loss of limbs or damage to internal organs. Extreme winter cold often causes poorly insulated water pipes to freeze. Even some poorly protected indoor plumbing may rupture as frozen water expands within them, causing property damage. Fires, paradoxically, become more hazardous during extreme cold. Water mains may break and water supplies may become unreliable, making firefighting more difficult. Cold waves that bring unexpected freezes and frosts during the growing season in mid-latitude zones can kill plants during the early and most vulnerable stages of growth. This results in crop failure as plants are killed before they can be harvested economically. Such cold waves have caused famines. Cold waves can also cause soil particles to harden and freeze, making it harder for plants and vegetation to grow within these areas. One extreme was the so-called Year Without a Summer of 1816, one of several years during the 1810s in which numerous crops failed during freakish summer cold snaps after volcanic eruptions reduced incoming sunlight. In some cases more frequent extremely cold winter weather – i.e. across parts of Asia and North America including the February 2021 North American cold wave – can be a result of climate change such as due to changes in the Arctic. However, conclusions that link climate change to cold waves are considered to still be controversial. The JRC PESETA IV project concluded in 2020 that overall climate change will result in a decline in the intensity and frequency of extreme cold spells, with milder winters reducing fatalities from extreme cold, even if individual cold extreme weather may sometimes be caused by changes due to climate change and possibly even become more frequent in some regions. According to a 2023 study, "weak extreme cold events (ECEs) significantly decrease in frequency, projection area and total area over the north hemisphere with global warming. However, the frequency, projection area and total area of strong ECEs show no significant trend, whereas they are increasing in Siberia and Canada." Heavy rain and storms Tropical cyclones Causes and attribution Attribution research Generally speaking, one event in extreme weather cannot be attributed to any one cause. However, certain system wide changes to global weather systems can lead to increased frequency or intensity of extreme weather events. Early research in extreme weather focused on statements about predicting certain events. Contemporary research focuses more on the attribution of causes to trends in events. In particular the field is focusing on climate change alongside other causal factors for these events. A 2016 report from the National Academies of Sciences, Engineering, and Medicine, recommended investing in improved shared practices across the field working on attribution research, improving the connection between research outcomes and weather forecasting. As more research is done in this area, scientists have begun to investigate the connection between climate change and extreme weather events and what future impacts may arise. Much of this work is done through climate modeling. Climate models provide important predictions about the future characteristics of the atmosphere, oceans, and Earth using data collected in the modern day. However, while climate models are vital for studying more complex processes such as climate change or ocean acidification, they are still only approximations. Moreover, weather events are complex and cannot be tied to a singular cause—there are often many atmospheric variables such as temperature, pressure, or moisture to note on top of any influences from climate change or natural variability. Natural variability Aspects of our climate system have a certain level of natural variability, and extreme weather events can occur for several reasons beyond human impact, including changes in pressure or the movement of air. Areas along the coast or located in tropical regions are more likely to experience storms with heavy precipitation than temperate regions, although such events can occur. The atmosphere is a complex and dynamic system, influenced by several factors such as the natural tilt and orbit of the Earth, the absorption or reflection of solar radiation, the movement of air masses, and the water cycle. Due to this, weather patterns can experience some variation, and so extreme weather can be attributed, at least in part, to the natural climate variability that exists on Earth. Climatic phenomena such as the El Niño-Southern Oscillation (ENSO) or the North Atlantic Oscillation (NAO) impact weather patterns in specific regions of the world, influencing temperature and precipitation. The record-breaking extreme weather events that have been catalogued throughout the past two hundred years most likely arise when climate patterns like ENSO or NAO work "in the same direction as human‐induced warming." Climate change Some studies assert a connection between rapidly warming arctic temperatures and thus a vanishing cryosphere to extreme weather in mid-latitudes. In a study published in Nature in 2019, scientists used several simulations to determine that the melting of ice sheets in Greenland and Antarctica could affect overall sea level and sea temperature. Other models have shown that modern temperature rise and the subsequent addition of meltwater to the ocean could lead to a disruption of the thermohaline circulation, which is responsible for the movement of seawater and distribution of heat around the globe. A collapse of this circulation in the northern hemisphere could lead to an increase in extreme temperatures in Europe, as well as more frequent storms by throwing off natural climate variability and conditions. Thus, as increasing temperatures cause glaciers to melt, mid-latitudes could experience shifts in weather patterns or temperatures. There were around 6,681 climate-related events reported during 2000-2019, compared to 3,656 climate-related events reported during 1980–1999. In this report, a 'climate-related event' refers to floods, storms, droughts, landslides, extreme temperatures (like heat waves or freezes), and wildfires; it excludes geophysical events such as volcanic eruptions, earthquakes, or mass movements. While there is evidence that a changing global climate, such as an increase in temperature, has impacted the frequency of extreme weather events, the most significant effects are likely to arise in the future. This is where climate models are useful, for they can provide simulations of how the atmosphere may behave over time and what steps need to be taken in the present day to mitigate any negative changes. The increasing probability of record week-long heat extremes occurrence depends on warming rate, rather than global warming level. Some researchers attribute increases in extreme weather occurrences to more reliable reporting systems. A difference in what qualifies as 'extreme weather' in varying climate systems could also be argued. Over or under reporting of casualties or losses can lead to inaccuracy in the impact of extreme weather. However, the UN reports show that, although some countries have experienced greater effects, there have been increases in extreme weather events on all continents. Current evidence and climate models show that an increasing global temperature will intensify extreme weather events around the globe, thereby amplifying human loss, damages and economic costs, and ecosystem destruction. Tropical cyclones and climate change In 2020, the National Oceanic and Atmospheric Administration (NOAA) of the U.S. government predicted that, over the 21st Century, the frequency of tropical storms and Atlantic hurricanes would decline by 25 percent while their maximum intensity would rise by 5 percent. Human activities that exacerbate the effects There are plenty of anthropogenic activities that can exacerbate the effects of extreme weather events. Urban planning often amplifies urban flooding impacts, especially in areas that are at increased risk of storms due to their location and climate variability. First, increasing the amount of impervious surfaces, such as sidewalks, roads, and roofs, means that less of the water from incoming storms is absorbed by the land. The destruction of wetlands, which act as a natural reservoir by absorbing water, can intensify the impact of floods and extreme precipitation. This can happen both inland and at the coast. However, wetland destruction along the coast can mean decreasing an area's natural 'cushion,' thus allowing storm surges and flood waters to reach farther inland during hurricanes or cyclones. Building homes below sea level or along a floodplain puts residents at increased risk of destruction or injury in an extreme precipitation event. More urban areas can also contribute to the rise of extreme or unusual weather events. Tall structures can alter the way that wind moves throughout an urban area, pushing warmer air upwards and inducing convection, creating thunderstorms. With these thunderstorms comes increased precipitation, which, because of the large amounts of impervious surfaces in cities, can have devastating impacts. Impervious surfaces also absorb energy from the sun and warm the atmosphere, causing drastic increases in temperatures in urban areas. This, along with pollution and heat released from cars and other anthropogenic sources, contributes to urban heat islands. Effects The effects of extreme weather includes, but are not limited to: Too much rain (heavy downpours), causing floods and landslides Too much heat and no rain (heatwave) causing droughts and wildfires Strong winds, such as hurricanes and tornadoes, causing damage to man made structures and animal habitats Large snowfalls, causing avalanches and blizzards Economic cost In the face of record breaking extreme weather events, climate change adaptation efforts fall short while economists are confronted with inflation, the cost-of-living crisis, and economic uncertainty. In 2011 the IPCC estimated, that annual losses have ranged since 1980 from a few billion to above US$200 billion, with the highest economic losses occurring in 2005, the year of Hurricane Katrina. The global weather-related disaster losses, such as loss of human lives, cultural heritage, and ecosystem services, are difficult to value and monetize, and thus they are poorly reflected in estimates of losses. The World Economic Forum Global Risks Perception Survey 2023-2024 (GRPS) found that 66 percent of respondents selected extreme weather as top risk. The survey was conducted after the 2023 heat waves. According to the GRPS results, the perception of necessary short and long-term risk management varies. Younger respondents prioritize environmental risks, including extreme weather, in the short-term. Respondents working in the private sector prioritize environmental risks as long-term. Loss of human lives The death toll from natural disasters has declined over 90 percent since the 1920s, according to the International Disaster Database, even as the total human population on Earth quadrupled, and temperatures rose 1.3 °C. In the 1920s, 5.4 million people died from natural disasters while in the 2010s, just 400,000 did. The most dramatic and rapid declines in deaths from extreme weather events have taken place in south Asia. Where a tropical cyclone in 1991 in Bangladesh killed 135,000 people, and a 1970 cyclone killed 300,000, the similarly-sized Cyclone Ampham, which struck India and Bangladesh in 2020, killed just 120 people in total. On July 23, 2020, Munich Re announced that the 2,900 total global deaths from natural disasters for the first half of 2020 were a record-low, and "much lower than the average figures for both the last 30 years and the last 10 years." A 2021 study found that 9.4% of global deaths between 2000 and 2019 – ~5 million annually – can be attributed to extreme temperature with cold-related ones making up the larger share and decreasing and heat-related ones making up ~0.91 % and increasing. Droughts and floods Climate change has led to an increase in the frequency and/or intensity of certain types of extreme weather. Storms such as hurricanes or tropical cyclones may experience greater rainfall, causing major flooding events or landslides by saturating soil. This is because warmer air is able to 'hold' more moisture due to the water molecules having increased kinetic energy, and precipitation occurs at a greater rate because more molecules have the critical speed needed to fall as rain drops. A shift in rainfall patterns can lead to greater amounts of precipitation in one area while another experiences much hotter, drier conditions, which can lead to drought. This is because an increase in temperatures also lead to an increase in evaporation at the surface of the earth, so more precipitation does not necessarily mean universally wetter conditions or a worldwide increase in drinking water.	Physical sciences	Climate change	Earth science
5759935	https://en.wikipedia.org/wiki/Thiamine%20deficiency	Thiamine deficiency	Thiamine deficiency is a medical condition of low levels of thiamine (vitamin B1). A severe and chronic form is known as beriberi. The name beriberi was possibly borrowed in the 18th century from the Sinhalese phrase (bæri bæri, “I cannot, I cannot”), owing to the weakness caused by the condition. The two main types in adults are wet beriberi and dry beriberi. Wet beriberi affects the cardiovascular system, resulting in a fast heart rate, shortness of breath, and leg swelling. Dry beriberi affects the nervous system, resulting in numbness of the hands and feet, confusion, trouble moving the legs, and pain. A form with loss of appetite and constipation may also occur. Another type, acute beriberi, found mostly in babies, presents with loss of appetite, vomiting, lactic acidosis, changes in heart rate, and enlargement of the heart. Risk factors include a diet of mostly white rice, alcoholism, dialysis, chronic diarrhea, and taking high doses of diuretics. In rare cases, it may be due to a genetic condition that results in difficulties absorbing thiamine found in food. Wernicke encephalopathy and Korsakoff syndrome are forms of dry beriberi. Diagnosis is based on symptoms, low levels of thiamine in the urine, high blood lactate, and improvement with thiamine supplementation. Treatment is by thiamine supplementation, either by mouth or by injection. With treatment, symptoms generally resolve in a few weeks. The disease may be prevented at the population level through the fortification of food. Thiamine deficiency is rare in the United States. It remains relatively common in sub-Saharan Africa. Outbreaks have been seen in refugee camps. Thiamine deficiency has been described for thousands of years in Asia, and became more common in the late 1800s with the increased processing of rice. Signs and symptoms Symptoms of beriberi include weight loss, emotional disturbances, impaired sensory perception, weakness and pain in the limbs, and periods of irregular heart rate. Edema (swelling of bodily tissues) is common. It may increase the amount of lactic acid and pyruvic acid within the blood. In advanced cases, the disease may cause high-output cardiac failure and death. Symptoms may occur concurrently with those of Wernicke's encephalopathy, a primarily neurological thiamine deficiency-related condition. Beriberi is divided into four categories. The first three are historical and the fourth, gastrointestinal beriberi, was recognized in 2004: Dry beriberi especially affects the peripheral nervous system. Wet beriberi especially affects the cardiovascular system and other bodily systems. Infantile beriberi affects the babies of malnourished mothers. Gastrointestinal beriberi affects the digestive system and other bodily systems. Dry beriberi Dry beriberi causes wasting and partial paralysis resulting from damaged peripheral nerves. It is also referred to as endemic neuritis. It is characterized by: Difficulty with walking Tingling or loss of sensation (numbness) in hands and feet Loss of tendon reflexes Loss of muscle function or paralysis of the lower legs Mental confusion/speech difficulties Pain Involuntary eye movements (nystagmus) Vomiting A selective impairment of the large proprioceptive sensory fibers without motor impairment can occur and present as a prominent sensory ataxia, which is a loss of balance and coordination due to loss of the proprioceptive inputs from the periphery and loss of position sense. Brain disease Wernicke's encephalopathy (WE), Korsakoff syndrome (also called alcohol amnestic disorder), and Wernicke–Korsakoff syndrome are forms of dry beriberi. Wernicke's encephalopathy is the most frequently encountered manifestation of thiamine deficiency in Western society, though it may also occur in patients with impaired nutrition from other causes, such as gastrointestinal disease, those with HIV/AIDS, and with the injudicious administration of parenteral glucose or hyperalimentation without adequate B-vitamin supplementation. This is a striking neuro-psychiatric disorder characterized by paralysis of eye movements, abnormal stance and gait, and markedly deranged mental function. Korsakoff syndrome, in general, is considered to occur with deterioration of brain function in patients initially diagnosed with WE. This is an amnestic-confabulatory syndrome characterized by retrograde and anterograde amnesia, impairment of conceptual functions, and decreased spontaneity and initiative. Alcoholics may have thiamine deficiency because of: Inadequate nutritional intake: Alcoholics tend to intake less than the recommended amount of thiamine. Decreased uptake of thiamine from the GI tract: Active transport of thiamine into enterocytes is disturbed during acute alcohol exposure. Liver thiamine stores are reduced due to hepatic steatosis or fibrosis. Impaired thiamine utilization: Magnesium, which is required for the binding of thiamine to thiamine-using enzymes within the cell, is also deficient due to chronic alcohol consumption. The inefficient use of any thiamine that does reach the cells will further exacerbate the thiamine deficiency. Ethanol per se inhibits thiamine transport in the gastrointestinal system and blocks phosphorylation of thiamine to its cofactor form (ThDP). Following improved nutrition and the removal of alcohol consumption, some impairments linked with thiamine deficiency are reversed, in particular poor brain functionality, although in more severe cases, Wernicke–Korsakoff syndrome leaves permanent damage. (See delirium tremens.) Wet beriberi Wet beriberi affects the heart and circulatory system. It is sometimes fatal, as it causes a combination of heart failure and weakening of the capillary walls, which causes the peripheral tissues to become edematous. Wet beriberi is characterized by: Increased heart rate Vasodilation leading to decreased systemic vascular resistance, and high-output heart failure Elevated jugular venous pressure Dyspnea (shortness of breath) on exertion Paroxysmal nocturnal dyspnea Peripheral edema (swelling of lower legs) or generalized edema (swelling throughout the body) Dilated cardiomyopathy Gastrointestinal beriberi Gastrointestinal beriberi causes abdominal pain. It is characterized by: Abdominal pain Nausea Vomiting Lactic acidosis Infants Infantile beriberi usually occurs between two and six months of age in children whose mothers have inadequate thiamine intake. It may present as either wet or dry beriberi. In the acute form, the baby develops dyspnea and cyanosis and soon dies of heart failure. These symptoms may be described in infantile beriberi: Hoarseness, where the child makes moves to moan, but emits no sound or just faint moans caused by nerve paralysis Weight loss, becoming thinner and then marasmic as the disease progresses Vomiting Diarrhea Pale skin Edema Ill temper Alterations of the cardiovascular system, especially tachycardia (rapid heart rate) Convulsions occasionally observed in the terminal stages Cause Beriberi is often caused by eating a diet with a very high proportion of calorie rich polished rice (common in Asia) or cassava root (common in sub-Saharan Africa), without much if any thiamine-containing animal products or vegetables. It may also be caused by shortcomings other than inadequate intake – diseases or operations on the digestive tract, alcoholism, dialysis or genetic deficiencies. All those causes mainly affect the central nervous system, and provoke the development of Wernicke's encephalopathy. Wernicke's disease is one of the most prevalent neurological or neuropsychiatric diseases. In autopsy series, features of Wernicke lesions are observed in approximately 2% of general cases. Medical record research shows that about 85% had not been diagnosed, although only 19% would be asymptomatic. In children, only 58% were diagnosed. In alcohol abusers, autopsy series showed neurological damages at rates of 12.5% or more. Mortality caused by Wernicke's disease reaches 17% of diseases, which means 3.4/1000 or about 25 million contemporaries. The number of people with Wernicke's disease may be even higher, considering that early stages may have dysfunctions prior to the production of observable lesions at necropsy. In addition, uncounted numbers of people can experience fetal damage and subsequent diseases. Genetics Genetic diseases of thiamine transport are rare but serious. Thiamine responsive megaloblastic anemia syndrome (TRMA) with diabetes mellitus and sensorineural deafness is an autosomal recessive disorder caused by mutations in the gene SLC19A2, a high affinity thiamine transporter. TRMA patients do not show signs of systemic thiamine deficiency, suggesting redundancy in the thiamine transport system. This has led to the discovery of a second high-affinity thiamine transporter, SLC19A3. Leigh disease (subacute necrotising encephalomyelopathy) is an inherited disorder that affects mostly infants in the first years of life and is invariably fatal. Pathological similarities between Leigh disease and WE led to the hypothesis that the cause was a defect in thiamine metabolism. One of the most consistent findings has been an abnormality of the activation of the pyruvate dehydrogenase complex. Mutations in the SLC19A3 gene have been linked to biotin-thiamine responsive basal ganglia disease, which is treated with pharmacological doses of thiamine and biotin, another B vitamin. Other disorders in which a putative role for thiamine has been implicated include subacute necrotising encephalomyelopathy, opsoclonus myoclonus syndrome (a paraneoplastic syndrome), and Nigerian seasonal ataxia (or African seasonal ataxia). In addition, several inherited disorders of ThDP-dependent enzymes have been reported, which may respond to thiamine treatment. Pathophysiology Thiamine in the human body has a half-life of 17 days and is quickly exhausted, particularly when metabolic demands exceed intake. A derivative of thiamine, thiamine pyrophosphate (TPP), is a cofactor involved in the citric acid cycle, as well as connecting the breakdown of sugars with the citric acid cycle. The citric acid cycle is a central metabolic pathway involved in the regulation of carbohydrate, lipid, and amino acid metabolism, and its disruption due to thiamine deficiency inhibits the production of many molecules including the neurotransmitters glutamic acid and GABA. Additionally, thiamine may also be directly involved in neuromodulation. Diagnosis A positive diagnosis test for thiamine deficiency involves measuring the activity of the enzyme transketolase in erythrocytes (Erythrocyte transketolase activation assay). Alternatively, thiamine and its phosphorylated derivatives can directly be detected in whole blood, tissues, foods, animal feed, and pharmaceutical preparations following the conversion of thiamine to fluorescent thiochrome derivatives (thiochrome assay) and separation by high-performance liquid chromatography (HPLC). Capillary electrophoresis (CE) techniques and in-capillary enzyme reaction methods have emerged as alternative techniques in quantifying and monitoring thiamine levels in samples. The normal thiamine concentration in EDTA-blood is about 20–100 μg/L. Treatment Many people with beriberi can be treated with thiamine alone. Given thiamine intravenously (and later orally), rapid and dramatic recovery occurs, generally within 24 hours. Improvements of peripheral neuropathy may require several months of thiamine treatment. Epidemiology Beriberi is a recurrent nutritional disease in detention houses, even in this century. In 1999, an outbreak of beriberi occurred in a detention center in Taiwan. High rates of illness and death from beriberi in overcrowded Haitian jails in 2007 were traced to the traditional practice of washing rice before cooking; this removed a nutritious coating which had been applied to the rice after processing (enriched white rice). In the Ivory Coast, among a group of prisoners with heavy punishment, 64% were affected by beriberi. Before beginning treatment, prisoners exhibited symptoms of dry or wet beriberi with neurological signs (tingling: 41%), cardiovascular signs (dyspnoea: 42%, thoracic pain: 35%), and edemas of the lower limbs (51%). With treatment, the rate of healing was about 97%. Populations under extreme stress may be at higher risk for beriberi. Displaced populations, such as refugees from war, are susceptible to micronutritional deficiency, including beriberi. The severe nutritional deprivation caused by famine also can cause beriberis, although symptoms may be overlooked in clinical assessment or masked by other famine-related problems. An extreme weight-loss diet can, rarely, induce a famine-like state and the accompanying beriberi. Workers on Chinese squid ships are at elevated risk of beriberi due to the simple carbohydrate-rich diet they are fed and the long period of time between shoring. Between 2013 and 2021, 15 workers on 14 ships have died with symptoms of beriberi. History Earliest written descriptions of thiamine deficiency are from ancient China in the context of Chinese medicine. One of the earliest is by Ge Hong in his book Zhou hou bei ji fang (Emergency Formulas to Keep up Your Sleeve) written sometime during the third century. Hong called the illness by the name jiao qi, which can be interpreted as "foot qi". He described the symptoms to include swelling, weakness, and numbness of the feet. He also acknowledged that the illness could be deadly, and claimed that it could be cured by eating certain foods, such as fermented soybeans in wine. Better known examples of early descriptions of "foot qi" are by Chao Yuanfang (who lived during 550–630) in his book Zhu bing yuan hou lun (Sources and Symptoms of All Diseases) and by Sun Simiao (581–682) in his book Bei ji qian jin yao fang (Essential Emergency Formulas Worth a Thousand in Gold). In the mid-19th century, interest in beriberi steadily rose as the disease became more noticeable with changes in diet in East and Southeast Asia. There was a steady uptick in medical publications, reaching one hundred and eighty-one publications from 1880 and 1889, and hundreds more in the following decades. The link to white rice was clear to Western doctors, but a confounding factor was that some other foods like meat failed to prevent beriberi, so it could not be easily explained as a lack of known chemicals like carbon or nitrogen. With no knowledge of vitamins, the etiology of beriberi was among the most hotly debated subjects in Victorian medicine. The first successful preventative measure against beriberi was discovered by Takaki Kanehiro, a British-trained Japanese medical doctor of the Imperial Japanese Navy, in the mid-1880s. Beriberi was a serious problem in the Japanese navy; sailors fell ill an average of four times a year in the period 1878 to 1881, and 35% were cases of beriberi. In 1882, Takaki learned of a very high incidence of beriberi among cadets on a training mission from Japan to Hawaii, via New Zealand and South America. The voyage lasted more than nine months and resulted in 169 cases of sickness and 25 deaths on a ship of 376 men. Takaki observed that beriberi was common among low-ranking crew who were often provided free rice, thus ate little else, but not among crews of Western navies, nor among Japanese officers who consumed a more varied diet. With the support of the Japanese Navy, he conducted an experiment in which another ship was deployed on the same route, except that its crew was fed a diet of meat, fish, barley, rice, and beans. At the end of the voyage, this crew had only 14 cases of beriberi and no deaths. This emphasis on varied diet contradicted observations by other doctors, and Takaki's carbon-based etiology was just as incorrect as similar theories before him, but the results of his experiment impressed the Japanese Navy, which adopted his proposed solution. By 1887 beriberi had been completely eliminated on Navy ships. In the same year, Takaki's experiment was described favorably in The Lancet, but his incorrect etiology was not taken seriously. In 1897, Christiaan Eijkman, a Dutch physician and pathologist, published his mid-1880s experiments showing that feeding unpolished rice (instead of the polished variety) to chickens helped to prevent beriberi. This was the first experiment to show that not a major chemical but some minor nutrient was the true cause of beriberi. The following year, Sir Frederick Hopkins postulated that some foods contained "accessory factors"—in addition to proteins, carbohydrates, fats, and salt—that were necessary for the functions of the human body. In 1901, Gerrit Grijns, a Dutch physician and assistant to Christiaan Eijkman in the Netherlands, correctly interpreted beriberi as a deficiency syndrome, and between 1910 and 1913, Edward Bright Vedder established that an extract of rice bran is a treatment for beriberi. In 1929, Eijkman and Hopkins were awarded the Nobel Prize for Physiology or Medicine for their discoveries. Japanese Army denialism Although the identification of beriberi as a deficiency syndrome was proven beyond a doubt by 1913, a Japanese group headed by Mori Ōgai and backed by Tokyo Imperial University continued to deny this conclusion until 1926. In 1886, Mori, then working in the Japanese Army Medical Bureau, asserted that white rice was sufficient as a diet for soldiers. Simultaneously, Navy surgeon general Takaki Kanehiro published the groundbreaking results described above. Mori, who had been educated under German doctors, responded that Takaki was a "fake doctor" due to his lack of prestigious medical background, while Mori himself and his fellow graduates of Tokyo Imperial University constituted the only "real doctors" in Japan and that they alone were capable of "experimental induction", although Mori himself had not conducted any beriberi experiments. The Japanese Navy sided with Takaki and adopted his suggestions. In order to prevent himself and the Army from losing face, Mori assembled a team of doctors and professors from Tokyo Imperial University and the Japanese Army who proposed that beriberi was caused by an unknown pathogen, which they described as etowasu (from the German Etwas, meaning "something"). They employed various social tactics to denounce vitamin deficiency experiments and prevent them from being published, while beriberi ravaged the Japanese Army. During the First Sino-Japanese War and Russo-Japanese War, Army soldiers continued to die in mass numbers from beriberi, while Navy sailors survived. In response to this severe loss of life, in 1907, the Army ordered the formation of a Beriberi Emergency Research Council, headed by Mori. Its members pledged to find the cause of beriberi. By 1919, with most Western doctors acknowledging that beriberi was a deficiency syndrome, the Emergency Research Council began conducting experiments using various vitamins, but stressed that "more research was necessary". During this period, more than 300,000 Japanese soldiers contracted beriberi and over 27,000 died. Mori died in 1922. The Beriberi Research Council disbanded in 1925, and by the time Eijkman and Hopkins were awarded the Nobel Prize, all of its members had acknowledged that beriberi was a deficiency syndrome. Etymology Although according to the Oxford English Dictionary, the term "beriberi" comes from a Sinhalese phrase meaning "weak, weak" or "I cannot, I cannot", the word being duplicated for emphasis, the origin of the phrase is questionable. It has also been suggested to come from Hindi, Arabic, and a few other languages, with many meanings like "weakness", "sailor", and even "sheep". Such suggested origins were listed by Heinrich Botho Scheube, among others. Edward Vedder wrote in his book Beriberi (1913) that "it is impossible to definitely trace the origin of the word beriberi". The word berbere was used in writing at least as early as 1568 by Diogo do Couto, when he described the deficiency in India. , which is a Japanese synonym for thiamine deficiency, comes from the way "jiao qi" is pronounced in Japanese. "Jiao qi is an old word used in Chinese medicine to describe beriberi. "Kakke is supposed to have entered into the Japanese language sometime between the sixth and eighth centuries. Other animals Poultry Mature chickens show signs three weeks after being fed a deficient diet. In young chicks, it can appear before two weeks of age. Onset is sudden in young chicks, with anorexia and an unsteady gait. Later on, locomotor signs begin, with an apparent paralysis of the flexor of the toes. The characteristic position is called "stargazing", with the affected animal sitting on its hocks with its head thrown back in a posture called opisthotonos. Response to administration of the vitamin is rather quick, occurring a few hours later. Ruminants Polioencephalomalacia (PEM) is the most common thiamine deficiency disorder in young ruminant and nonruminant animals. Symptoms of PEM include a profuse, but transient, diarrhea, listlessness, circling movements, stargazing or opisthotonus (head drawn back over neck), and muscle tremors. The most common cause is high-carbohydrate feeds, leading to the overgrowth of thiaminase-producing bacteria, but dietary ingestion of thiaminase (e.g., in bracken fern), or inhibition of thiamine absorption by high sulfur intake are also possible. Another cause of PEM is Clostridium sporogenes or Bacillus aneurinolyticus infection. These bacteria produce thiaminases that can cause an acute thiamine deficiency in the affected animal. Snakes Snakes that consume a diet largely composed of goldfish and feeder minnows are susceptible to developing thiamine deficiency. This is often a problem observed in captivity when keeping garter and ribbon snakes that are fed a goldfish-exclusive diet, as these fish contain thiaminase, an enzyme that breaks down thiamine. Wild birds and fish Thiamine deficiency has been identified as the cause of a paralytic disease affecting wild birds in the Baltic Sea area dating back to 1982. In this condition, there is difficulty in keeping the wings folded along the side of the body when resting, loss of the ability to fly and voice, with eventual paralysis of the wings and legs and death. It affects primarily 0.5–1 kg-sized birds such as the European herring gull (Larus argentatus), common starling (Sturnus vulgaris), and common eider (Somateria mollissima). Researchers noted, "Because the investigated species occupy a wide range of ecological niches and positions in the food web, we are open to the possibility that other animal classes may develop thiamine deficiency, as well."p. 12006 In the counties of Blekinge and Skåne, mass deaths of several bird species, especially the European herring gull, have been observed since the early 2000s. More recently, species of other classes seems to be affected. High mortality of salmon (Salmo salar) in the river Mörrumsån is reported, and mammals such as the Eurasian elk (Alces alces) have died in unusually high numbers. Lack of thiamine is the common denominator where analysis is done. In April 2012, the County Administrative Board of Blekinge found the situation so alarming that they asked the Swedish government to set up a closer investigation.	Biology and health sciences	Health and fitness: General	Health
187315	https://en.wikipedia.org/wiki/Antenna%20%28zoology%29	Antenna (zoology)	Antennae (: antenna) (sometimes referred to as "feelers") are paired appendages used for sensing in arthropods. Antennae are connected to the first one or two segments of the arthropod head. They vary widely in form but are always made of one or more jointed segments. While they are typically sensory organs, the exact nature of what they sense and how they sense it is not the same in all groups. Functions may variously include sensing touch, air motion, heat, vibration (sound), and especially smell or taste. Antennae are sometimes modified for other purposes, such as mating, brooding, swimming, and even anchoring the arthropod to a substrate. Larval arthropods have antennae that differ from those of the adult. Many crustaceans, for example, have free-swimming larvae that use their antennae for swimming. Antennae can also locate other group members if the insect lives in a group, like the ant. The common ancestor of all arthropods likely had one pair of uniramous (unbranched) antenna-like structures, followed by one or more pairs of biramous (having two major branches) leg-like structures, as seen in some modern crustaceans and fossil trilobites. Except for the chelicerates and proturans, which have none, all non-crustacean arthropods have a single pair of antennae. Crustaceans Crustaceans bear two pairs of antennae. The pair attached to the first segment of the head are called primary antennae or antennules. This pair is generally uniramous, but is biramous in crabs and lobsters and remipedes. The pair attached to the second segment are called secondary antennae or simply antennae. The second antennae are plesiomorphically biramous, but many species later evolved uniramous pairs. The second antennae may be significantly reduced (e.g. remipedes) or apparently absent (e.g. barnacles). The subdivisions of crustacean antennae have many names, including flagellomeres (a shared term with insects), annuli, articles, and segments. The terminal ends of crustacean antennae have two major categorizations: segmented and flagellate. An antenna is considered segmented if each of the annuli is separate from those around it and has individual muscle attachments. Flagellate antennae, on the other hand, have muscle attachments only around the base, acting as a hinge for the flagellum—a flexible string of annuli with no muscle attachment. There are several notable non-sensory uses of antennae in crustaceans. Many crustaceans have a mobile larval stage called a nauplius, which is characterized by its use of antennae for swimming. Barnacles, a highly modified crustacean, use their antennae to attach to rocks and other surfaces. The second antennae in the burrowing Hippoidea and Corystidae have setae that interlock to form a tube or "snorkel" which funnels filtered water over the gills. Insects Some claim insects evolved from prehistoric crustaceans, and they have secondary antennae like crustaceans, but not primary antennae. Antennae are the primary olfactory sensors of insects and are accordingly well-equipped with a wide variety of sensilla (singular: sensillum). Paired, mobile, and segmented, they are located between the eyes on the forehead. Embryologically, they represent the appendages of the second head segment. All insects have antennae, however they may be greatly reduced in the larval forms. Amongst the non-insect classes of the Hexapoda, both Collembola and Diplura have antenna, but Protura do not. Antennal fibrillae play an important role in Culex pipiens mating practices. The erection of these fibrillae is considered to be the first stage in reproduction. These fibrillae serve different functions across the sexes. As antennal fibrillae are used by female C. pipiens to locate hosts to feed on, male C. pipiens utilize them to locate female mates. Structure The three basic segments of the typical insect antenna are the scape or scapus (base), the pedicel or pedicellus (stem), and finally the flagellum, which often comprises many units known as flagellomeres. The pedicel (the second segment) contains the Johnston's organ which is a collection of sensory cells. The scape is mounted in a socket in a more or less ring-shaped sclerotised region called the torulus, often a raised portion of the insect's head capsule. The socket is closed off by the membrane into which the base of the scape is set. However, the antenna does not hang free on the membrane, but pivots on a rigidly sprung projection from the rim of the torulus. That projection on which the antenna pivots is called the antennifer. The whole structure enables the insect to move the antenna as a whole by applying internal muscles connected to the scape. The pedicel is flexibly connected to the distal end of the scape and its movements in turn can be controlled by muscular connections between the scape and pedicel. The number of flagellomeres can vary greatly between insect species, and often is of diagnostic importance. True flagellomeres are connected by membranous linkage that permits movement, though the flagellum of "true" insects does not have any intrinsic muscles. Some other Arthropoda do however have intrinsic muscles throughout the flagellum. Such groups include the Symphyla, Collembola and Diplura. In many true insects, especially the more primitive groups such as Thysanura and Blattodea, the flagellum partly or entirely consists of a flexibly connected string of small ring-shaped annuli. The annuli are not true flagellomeres, and in a given insect species the number of annuli generally is not as consistent as the number of flagellomeres in most species. In many beetles and in the chalcidoid wasps, the apical flagellomeres form a club shape, called the clava. The collective term for the segments between the club and the antennal base is the funicle; traditionally in describing beetle anatomy, the term "funicle" refers to the segments between the club and the scape. However, traditionally in working on wasps the funicle is taken to comprise the segments between the club and the pedicel. Quite commonly the funicle beyond the pedicel is quite complex in Endopterygota such as beetles, moths and Hymenoptera, and one common adaptation is the ability to fold the antenna in the middle, at the joint between the pedicel and the flagellum. This gives an effect like a "knee bend", and such an antenna is said to be geniculate. Geniculate antennae are common in the Coleoptera and Hymenoptera. They are important for insects like ants that follow scent trails, for bees and wasps that need to "sniff" the flowers that they visit, and for beetles such as Scarabaeidae and Curculionidae that need to fold their antennae away when they self-protectively fold up all their limbs in defensive attitudes. Because the funicle is without intrinsic muscles, it generally must move as a unit, in spite of being articulated. However, some funicles are complex and very mobile. For example, the Scarabaeidae have lamellate antennae that can be folded tightly for safety or spread openly for detecting odours or pheromones. The insect manages such actions by changes in blood pressure, by which it exploits elasticity in walls and membranes in the funicles, which are in effect erectile. In the groups with more uniform antennae (for example: millipedes), all segments are called antennomeres. Some groups have a simple or variously modified apical or subapical bristle called an arista (this may be especially well-developed in various Diptera). Functions Olfactory receptors on the antennae bind to free-floating molecules, such as water vapour, and odours including pheromones. The neurons that possess these receptors signal this binding by sending action potentials down their axons to the antennal lobe in the brain. From there, neurons in the antennal lobes connect to mushroom bodies that identify the odour. The sum of the electrical potentials of the antennae to a given odour can be measured using an electroantennogram. In the monarch butterfly, antennae are necessary for proper time-compensated solar compass orientation during migration. Antennal clocks exist in monarchs, and they are likely to provide the primary timing mechanism for sun compass orientation. In the African cotton leafworm, antennae have an important function in signaling courtship. Specifically, antennae are required for males to answer the female mating call. Although females do not require antennae for mating, a mating that resulted from a female without antennae was abnormal. In the diamondback moth, antennae serve to gather information about a host plant's taste and odor. After the desired taste and odor has been identified, the female moth will deposit her eggs onto the plant. Giant swallowtail butterflies also rely on antenna sensitivity to volatile compounds to identify host plants. It was found that females are actually more responsive with their antenna sensing, most likely because they are responsible for oviposition on the correct plant. In the crepuscular hawk moth (Manduca sexta), antennae aid in flight stabilization. Similar to halteres in Dipteran insects, the antennae transmit coriolis forces through the Johnston's organ that can then be used for corrective behavior. A series of low-light, flight stability studies in which moths with flagellae amputated near the pedicel showed significantly decreased flight stability over those with intact antennae. To determine whether there may be other antennal sensory inputs, a second group of moths had their antennae amputated and then re-attached, before being tested in the same stability study. These moths showed slightly decreased performance from intact moths, indicating there are possibly other sensory inputs used in flight stabilization. Re-amputation of the antennae caused a drastic decrease in flight stability to match that of the first amputated group.	Biology and health sciences	Nervous system	Biology
187317	https://en.wikipedia.org/wiki/Antenna%20%28radio%29	Antenna (radio)	In radio engineering, an antenna (American English) or aerial (British English) is an electronic device that converts an alternating electric current into radio waves (transmitting), or radio waves into an electric current (receiving). It is the interface between radio waves propagating through space and electric currents moving in metal conductors, used with a transmitter or receiver. In transmission, a radio transmitter supplies an electric current to the antenna's terminals, and the antenna radiates the energy from the current as electromagnetic waves (radio waves). In reception, an antenna intercepts some of the power of a radio wave in order to produce an electric current at its terminals, that is applied to a receiver to be amplified. Antennas are essential components of all radio equipment. An antenna is an array of conductors (elements), electrically connected to the receiver or transmitter. Antennas can be designed to transmit and receive radio waves in all horizontal directions equally (omnidirectional antennas), or preferentially in a particular direction (directional, or high-gain, or "beam" antennas). An antenna may include components not connected to the transmitter, parabolic reflectors, horns, or parasitic elements, which serve to direct the radio waves into a beam or other desired radiation pattern. Strong directivity and good efficiency when transmitting are hard to achieve with antennas with dimensions that are much smaller than a half wavelength. The first antennas were built in 1886 by German physicist Heinrich Hertz in his pioneering experiments to prove the existence of electromagnetic waves predicted by the 1867 electromagnetic theory of James Clerk Maxwell. Hertz placed dipole antennas at the focal point of parabolic reflectors for both transmitting and receiving. Starting in 1895, Guglielmo Marconi began development of antennas practical for long-distance, wireless telegraphy, for which he received the 1909 Nobel Prize in physics. Terminology The words antenna and aerial are used interchangeably. Occasionally the equivalent term "aerial" is used to specifically mean an elevated horizontal wire antenna. The origin of the word antenna relative to wireless apparatus is attributed to Italian radio pioneer Guglielmo Marconi. In the summer of 1895, Marconi began testing his wireless system outdoors on his father's estate near Bologna and soon began to experiment with long wire "aerials" suspended from a pole. In Italian a tent pole is known as l'antenna centrale, and the pole with the wire was simply called l'antenna. Until then wireless radiating transmitting and receiving elements were known simply as "terminals". Because of his prominence, Marconi's use of the word antenna spread among wireless researchers and enthusiasts, and later to the general public. Antenna may refer broadly to an entire assembly including support structure, enclosure (if any), etc., in addition to the actual RF current-carrying components. A receiving antenna may include not only the passive metal receiving elements, but also an integrated preamplifier or mixer, especially at and above microwave frequencies. Overview Antennas are required by any radio receiver or transmitter to couple its electrical connection to the electromagnetic field. Radio waves are electromagnetic waves which carry signals through the air (or through space) at the speed of light with almost no transmission loss. Antennas can be classified as omnidirectional, radiating energy approximately equally in all horizontal directions, or directional, where radio waves are concentrated in some direction(s). A so-called beam antenna is unidirectional, designed for maximum response in the direction of the other station, whereas many other antennas are intended to accommodate stations in various directions but are not truly omnidirectional. Since antennas obey reciprocity the same radiation pattern applies to transmission as well as reception of radio waves. A hypothetical antenna that radiates equally in all directions (vertical as well as all horizontal angles) is called an isotropic radiator; however, these cannot exist in practice nor would they be particularly desired. For most terrestrial communications, rather, there is an advantage in reducing radiation toward the sky or ground in favor of horizontal direction(s). A dipole antenna oriented horizontally sends no energy in the direction of the conductor – this is called the antenna null – but is usable in most other directions. A number of such dipole elements can be combined into an antenna array such as the Yagi–Uda in order to favor a single horizontal direction, thus termed a beam antenna. The dipole antenna, which is the basis for most antenna designs, is a balanced component, with equal but opposite voltages and currents applied at its two terminals. The vertical antenna is a monopole antenna, not balanced with respect to ground. The ground (or any large conductive surface) plays the role of the second conductor of a monopole. Since monopole antennas rely on a conductive surface, they may be mounted with a ground plane to approximate the effect of being mounted on the Earth's surface. More complex antennas increase the directivity of the antenna. Additional elements in the antenna structure, which need not be directly connected to the receiver or transmitter, increase its directionality. Antenna "gain" describes the concentration of radiated power into a particular solid angle of space. "Gain" is perhaps an unfortunately chosen term, by comparison with amplifier "gain" which implies a net increase in power. In contrast, for antenna "gain", the power increased in the desired direction is at the expense of power reduced in undesired directions. Unlike amplifiers, antennas are electrically "passive" devices which conserve total power, and there is no increase in total power above that delivered from the power source (the transmitter), only improved distribution of that fixed total. A phased array consists of two or more simple antennas which are connected together through an electrical network. This often involves a number of parallel dipole antennas with a certain spacing. Depending on the relative phase introduced by the network, the same combination of dipole antennas can operate as a "broadside array" (directional normal to a line connecting the elements) or as an "end-fire array" (directional along the line connecting the elements). Antenna arrays may employ any basic (omnidirectional or weakly directional) antenna type, such as dipole, loop or slot antennas. These elements are often identical. Log-periodic and frequency-independent antennas employ self-similarity in order to be operational over a wide range of bandwidths. The most familiar example is the log-periodic dipole array which can be seen as a number (typically 10 to 20) of connected dipole elements with progressive lengths in an endfire array making it rather directional; it finds use especially as a rooftop antenna for television reception. On the other hand, a Yagi–Uda antenna (or simply "Yagi"), with a somewhat similar appearance, has only one dipole element with an electrical connection; the other parasitic elements interact with the electromagnetic field in order to realize a highly directional antenna but with a narrow bandwidth. Even greater directionality can be obtained using aperture antennas such as the parabolic reflector or horn antenna. Since high directivity in an antenna depends on it being large compared to the wavelength, highly directional antennas (thus with high antenna gain) become more practical at higher frequencies (UHF and above). At low frequencies (such as AM broadcast), arrays of vertical towers are used to achieve directionality and they will occupy large areas of land. For reception, a long Beverage antenna can have significant directivity. For non directional portable use, a short vertical antenna or small loop antenna works well, with the main design challenge being that of impedance matching. With a vertical antenna a loading coil at the base of the antenna may be employed to cancel the reactive component of impedance; small loop antennas are tuned with parallel capacitors for this purpose. An antenna lead-in is the transmission line, or feed line, which connects the antenna to a transmitter or receiver. The "antenna feed" may refer to all components connecting the antenna to the transmitter or receiver, such as an impedance matching network in addition to the transmission line. In a so-called "aperture antenna", such as a horn or parabolic dish, the "feed" may also refer to a basic radiating antenna embedded in the entire system of reflecting elements (normally at the focus of the parabolic dish or at the throat of a horn) which could be considered the one active element in that antenna system. A microwave antenna may also be fed directly from a waveguide in place of a (conductive) transmission line. An antenna counterpoise, or ground plane, is a structure of conductive material which improves or substitutes for the ground. It may be connected to or insulated from the natural ground. In a monopole antenna, this aids in the function of the natural ground, particularly where variations (or limitations) of the characteristics of the natural ground interfere with its proper function. Such a structure is normally connected to the return connection of an unbalanced transmission line such as the shield of a coaxial cable. An electromagnetic wave refractor in some aperture antennas is a component which due to its shape and position functions to selectively delay or advance portions of the electromagnetic wavefront passing through it. The refractor alters the spatial characteristics of the wave on one side relative to the other side. It can, for instance, bring the wave to a focus or alter the wave front in other ways, generally in order to maximize the directivity of the antenna system. This is the radio equivalent of an optical lens. An antenna coupling network is a passive network (generally a combination of inductive and capacitive circuit elements) used for impedance matching in between the antenna and the transmitter or receiver. This may be used to minimize losses on the feed line, by reducing transmission line's standing wave ratio, and to present the transmitter or receiver with a standard resistive impedance needed for its optimum operation. The feed point location(s) is selected, and antenna elements electrically similar to tuner components may be incorporated in the antenna structure itself, to improve the match. Reciprocity It is a fundamental property of antennas that most of the electrical characteristics of an antenna, such as those described in the next section (e.g. gain, radiation pattern, impedance, bandwidth, resonant frequency and polarization), are the same whether the antenna is transmitting or receiving. For example, the "receiving pattern" (sensitivity to incoming signals as a function of direction) of an antenna when used for reception is identical to the radiation pattern of the antenna when it is driven and functions as a radiator, even though the current and voltage distributions on the antenna itself are different for receiving and sending. This is a consequence of the reciprocity theorem of electromagnetics. Therefore, in discussions of antenna properties no distinction is usually made between receiving and transmitting terminology, and the antenna can be viewed as either transmitting or receiving, whichever is more convenient. A necessary condition for the aforementioned reciprocity property is that the materials in the antenna and transmission medium are linear and reciprocal. Reciprocal (or bilateral) means that the material has the same response to an electric current or magnetic field in one direction, as it has to the field or current in the opposite direction. Most materials used in antennas meet these conditions, but some microwave antennas use high-tech components such as isolators and circulators, made of nonreciprocal materials such as ferrite. These can be used to give the antenna a different behavior on receiving than it has on transmitting, which can be useful in applications like radar. Resonant antennas The majority of antenna designs are based on the resonance principle. This relies on the behaviour of moving electrons, which reflect off surfaces where the dielectric constant changes, in a fashion similar to the way light reflects when optical properties change. In these designs, the reflective surface is created by the end of a conductor, normally a thin metal wire or rod, which in the simplest case has a feed point at one end where it is connected to a transmission line. The conductor, or element, is aligned with the electrical field of the desired signal, normally meaning it is perpendicular to the line from the antenna to the source (or receiver in the case of a broadcast antenna). The radio signal's electrical component induces a voltage in the conductor. This causes an electrical current to begin flowing in the direction of the signal's instantaneous field. When the resulting current reaches the end of the conductor, it reflects, which is equivalent to a 180 degree change in phase. If the conductor is of a wavelength long, current from the feed point will undergo 90 degree phase change by the time it reaches the end of the conductor, reflect through 180 degrees, and then another 90 degrees as it travels back. That means it has undergone a total 360 degree phase change, returning it to the original signal. The current in the element thus adds to the current being created from the source at that instant. This process creates a standing wave in the conductor, with the maximum current at the feed. The ordinary half-wave dipole is probably the most widely used antenna design. This consists of two wavelength elements arranged end-to-end, and lying along essentially the same axis (or collinear), each feeding one side of a two-conductor transmission wire. The physical arrangement of the two elements places them 180 degrees out of phase, which means that at any given instant one of the elements is driving current into the transmission line while the other is pulling it out. The monopole antenna is essentially one half of the half-wave dipole, a single wavelength element with the other side connected to ground or an equivalent ground plane (or counterpoise). Monopoles, which are one-half the size of a dipole, are common for long-wavelength radio signals where a dipole would be impractically large. Another common design is the folded dipole which consists of two (or more) half-wave dipoles placed side by side and connected at their ends but only one of which is driven. The standing wave forms with this desired pattern at the design operating frequency, , and antennas are normally designed to be this size. However, feeding that element with 3 (whose wavelength is that of ) will also lead to a standing wave pattern. Thus, an antenna element is also resonant when its length is of a wavelength. This is true for all odd multiples of wavelength. This allows some flexibility of design in terms of antenna lengths and feed points. Antennas used in such a fashion are known to be harmonically operated. Resonant antennas usually use a linear conductor (or element), or pair of such elements, each of which is about a quarter of the wavelength in length (an odd multiple of quarter wavelengths will also be resonant). Antennas that are required to be small compared to the wavelength sacrifice efficiency and cannot be very directional. Since wavelengths are so small at higher frequencies (UHF, microwaves) trading off performance to obtain a smaller physical size is usually not required. Current and voltage distribution The quarter-wave elements imitate a series-resonant electrical element due to the standing wave present along the conductor. At the resonant frequency, the standing wave has a current peak and voltage node (minimum) at the feed. In electrical terms, this means that at that position, the element has minimum impedance magnitude, generating the maximum current for minimum voltage. This is the ideal situation, because it produces the maximum output for the minimum input, producing the highest possible efficiency. Contrary to an ideal (lossless) series-resonant circuit, a finite resistance remains (corresponding to the relatively small voltage at the feed-point) due to the antenna's resistance to radiating, as well as any conventional electrical losses from producing heat. Recall that a current will reflect when there are changes in the electrical properties of the material. In order to efficiently transfer the received signal into the transmission line, it is important that the transmission line has the same impedance as its connection point on the antenna, otherwise some of the signal will be reflected backwards into the body of the antenna; likewise part of the transmitter's signal power will be reflected back to transmitter, if there is a change in electrical impedance where the feedline joins the antenna. This leads to the concept of impedance matching, the design of the overall system of antenna and transmission line so the impedance is as close as possible, thereby reducing these losses. Impedance matching is accomplished by a circuit called an antenna tuner or impedance matching network between the transmitter and antenna. The impedance match between the feedline and antenna is measured by a parameter called the standing wave ratio (SWR) on the feedline. Consider a half-wave dipole designed to work with signals with wavelength 1 m, meaning the antenna would be approximately 50 cm from tip to tip. If the element has a length-to-diameter ratio of 1000, it will have an inherent impedance of about 63 ohms resistive. Using the appropriate transmission wire or balun, we match that resistance to ensure minimum signal reflection. Feeding that antenna with a current of 1 Ampere will require 63 Volts, and the antenna will radiate 63 Watts (ignoring losses) of radio frequency power. Now consider the case when the antenna is fed a signal with a wavelength of 1.25 m; in this case the current induced by the signal would arrive at the antenna's feedpoint out-of-phase with the signal, causing the net current to drop while the voltage remains the same. Electrically this appears to be a very high impedance. The antenna and transmission line no longer have the same impedance, and the signal will be reflected back into the antenna, reducing output. This could be addressed by changing the matching system between the antenna and transmission line, but that solution only works well at the new design frequency. The result is that the resonant antenna will efficiently feed a signal into the transmission line only when the source signal's frequency is close to that of the design frequency of the antenna, or one of the resonant multiples. This makes resonant antenna designs inherently narrow-band: Only useful for a small range of frequencies centered around the resonance(s). Electrically short antennas It is possible to use simple impedance matching techniques to allow the use of monopole or dipole antennas substantially shorter than the or wave, respectively, at which they are resonant. As these antennas are made shorter (for a given frequency) their impedance becomes dominated by a series capacitive (negative) reactance; by adding an appropriate size "loading coil" – a series inductance with equal and opposite (positive) reactance – the antenna's capacitive reactance may be cancelled leaving only a pure resistance. Sometimes the resulting (lower) electrical resonant frequency of such a system (antenna plus matching network) is described using the concept of electrical length, so an antenna used at a lower frequency than its resonant frequency is called an electrically short antenna For example, at 30 MHz (10 m wavelength) a true resonant wave monopole would be almost 2.5 meters long, and using an antenna only 1.5 meters tall would require the addition of a loading coil. Then it may be said that the coil has lengthened the antenna to achieve an electrical length of 2.5 meters. However, the resulting resistive impedance achieved will be quite a bit lower than that of a true wave (resonant) monopole, often requiring further impedance matching (a transformer) to the desired transmission line. For ever shorter antennas (requiring greater "electrical lengthening") the radiation resistance plummets (approximately according to the square of the antenna length), so that the mismatch due to a net reactance away from the electrical resonance worsens. Or one could as well say that the equivalent resonant circuit of the antenna system has a higher Q factor and thus a reduced bandwidth, which can even become inadequate for the transmitted signal's spectrum. Resistive losses due to the loading coil, relative to the decreased radiation resistance, entail a reduced electrical efficiency, which can be of great concern for a transmitting antenna, but bandwidth is the major factor that sets the size of antennas at 1 MHz and lower frequencies. Arrays and reflectors The radiant flux as a function of the distance from the transmitting antenna varies according to the inverse-square law, since that describes the geometrical divergence of the transmitted wave. For a given incoming flux, the power acquired by a receiving antenna is proportional to its effective area. This parameter compares the amount of power captured by a receiving antenna in comparison to the flux of an incoming wave (measured in terms of the signal's power density in watts per square metre). A half-wave dipole has an effective area of about 0.13 seen from the broadside direction. If higher gain is needed one cannot simply make the antenna larger. Due to the constraint on the effective area of a receiving antenna detailed below, one sees that for an already-efficient antenna design, the only way to increase gain (effective area) is by reducing the antenna's gain in another direction. If a half-wave dipole is not connected to an external circuit but rather shorted out at the feedpoint, then it becomes a resonant half-wave element which efficiently produces a standing wave in response to an impinging radio wave. Because there is no load to absorb that power, it retransmits all of that power, possibly with a phase shift which is critically dependent on the element's exact length. Thus such a conductor can be arranged in order to transmit a second copy of a transmitter's signal in order to affect the radiation pattern (and feedpoint impedance) of the element electrically connected to the transmitter. Antenna elements used in this way are known as passive radiators. A Yagi–Uda array uses passive elements to greatly increase gain in one direction (at the expense of other directions). A number of parallel approximately half-wave elements (of very specific lengths) are situated parallel to each other, at specific positions, along a boom; the boom is only for support and not involved electrically. Only one of the elements is electrically connected to the transmitter or receiver, while the remaining elements are passive. The Yagi produces a fairly large gain (depending on the number of passive elements) and is widely used as a directional antenna with an antenna rotor to control the direction of its beam. It suffers from having a rather limited bandwidth, restricting its use to certain applications. Rather than using one driven antenna element along with passive radiators, one can build an array antenna in which multiple elements are all driven by the transmitter through a system of power splitters and transmission lines in relative phases so as to concentrate the RF power in a single direction. What's more, a phased array can be made "steerable", that is, by changing the phases applied to each element the radiation pattern can be shifted without physically moving the antenna elements. Another common array antenna is the log-periodic dipole array which has an appearance similar to the Yagi (with a number of parallel elements along a boom) but is totally dissimilar in operation as all elements are connected electrically to the adjacent element with a phase reversal; using the log-periodic principle it obtains the unique property of maintaining its performance characteristics (gain and impedance) over a very large bandwidth. When a radio wave hits a large conducting sheet it is reflected (with the phase of the electric field reversed) just as a mirror reflects light. Placing such a reflector behind an otherwise non-directional antenna will insure that the power that would have gone in its direction is redirected toward the desired direction, increasing the antenna's gain by a factor of at least 2. Likewise, a corner reflector can insure that all of the antenna's power is concentrated in only one quadrant of space (or less) with a consequent increase in gain. Practically speaking, the reflector need not be a solid metal sheet, but can consist of a curtain of rods aligned with the antenna's polarization; this greatly reduces the reflector's weight and wind load. Specular reflection of radio waves is also employed in a parabolic reflector antenna, in which a curved reflecting surface effects focussing of an incoming wave toward a so-called feed antenna; this results in an antenna system with an effective area comparable to the size of the reflector itself. Other concepts from geometrical optics are also employed in antenna technology, such as with the lens antenna. Characteristics The antenna's power gain (or simply "gain") also takes into account the antenna's efficiency, and is often the primary figure of merit. Antennas are characterized by a number of performance measures which a user would be concerned with in selecting or designing an antenna for a particular application. A plot of the directional characteristics in the space surrounding the antenna is its radiation pattern. Bandwidth The frequency range or bandwidth over which an antenna functions well can be very wide (as in a log-periodic antenna) or narrow (as in a small loop antenna); outside this range the antenna impedance becomes a poor match to the transmission line and transmitter (or receiver). Use of the antenna well away from its design frequency affects its radiation pattern, reducing its directive gain. Generally an antenna will not have a feed-point impedance that matches that of a transmission line; a matching network between antenna terminals and the transmission line will improve power transfer to the antenna. A non-adjustable matching network will most likely place further limits the usable bandwidth of the antenna system. It may be desirable to use tubular elements, instead of thin wires, to make an antenna; these will allow a greater bandwidth. Or, several thin wires can be grouped in a cage to simulate a thicker element. This widens the bandwidth of the resonance. Amateur radio antennas that operate at several frequency bands which are widely separated from each other may connect elements resonant at those different frequencies in parallel. Most of the transmitter's power will flow into the resonant element while the others present a high impedance. Another solution uses traps, parallel resonant circuits which are strategically placed in breaks created in long antenna elements. When used at the trap's particular resonant frequency the trap presents a very high impedance (parallel resonance) effectively truncating the element at the location of the trap; if positioned correctly, the truncated element makes a proper resonant antenna at the trap frequency. At substantially higher or lower frequencies the trap allows the full length of the broken element to be employed, but with a resonant frequency shifted by the net reactance added by the trap. The bandwidth characteristics of a resonant antenna element can be characterized according to its where the resistance involved is the radiation resistance, which represents the emission of energy from the resonant antenna to free space. The of a narrow band antenna can be as high as 15. On the other hand, the reactance at the same off-resonant frequency of one using thick elements is much less, consequently resulting in a as low as 5. These two antennas may perform equivalently at the resonant frequency, but the second antenna will perform over a bandwidth 3 times as wide as the antenna consisting of a thin conductor. Antennas for use over much broader frequency ranges are achieved using further techniques. Adjustment of a matching network can, in principle, allow for any antenna to be matched at any frequency. Thus the small loop antenna built into most AM broadcast (medium wave) receivers has a very narrow bandwidth, but is tuned using a parallel capacitance which is adjusted according to the receiver tuning. On the other hand, log-periodic antennas are not resonant at any single frequency but can (in principle) be built to attain similar characteristics (including feedpoint impedance) over any frequency range. These are therefore commonly used (in the form of directional log-periodic dipole arrays) as television antennas. Gain Gain is a parameter which measures the degree of directivity of the antenna's radiation pattern. A high-gain antenna will radiate most of its power in a particular direction, while a low-gain antenna will radiate over a wide angle. The antenna gain, or power gain of an antenna is defined as the ratio of the intensity (power per unit surface area) radiated by the antenna in the direction of its maximum output, at an arbitrary distance, divided by the intensity radiated at the same distance by a hypothetical isotropic antenna which radiates equal power in all directions. This dimensionless ratio is usually expressed logarithmically in decibels, these units are called decibels-isotropic (dBi) A second unit used to measure gain is the ratio of the power radiated by the antenna to the power radiated by a half-wave dipole antenna ; these units are called decibels-dipole (dBd) Since the gain of a half-wave dipole is 2.15 dBi and the logarithm of a product is additive, the gain in dBi is just 2.15 decibels greater than the gain in dBd High-gain antennas have the advantage of longer range and better signal quality, but must be aimed carefully at the other antenna. An example of a high-gain antenna is a parabolic dish such as a satellite television antenna. Low-gain antennas have shorter range, but the orientation of the antenna is relatively unimportant. An example of a low-gain antenna is the whip antenna found on portable radios and cordless phones. Antenna gain should not be confused with amplifier gain, a separate parameter measuring the increase in signal power due to an amplifying device placed at the front-end of the system, such as a low-noise amplifier. Effective area or aperture The effective area or effective aperture of a receiving antenna expresses the portion of the power of a passing electromagnetic wave which the antenna delivers to its terminals, expressed in terms of an equivalent area. For instance, if a radio wave passing a given location has a flux of 1 pW / m2 (10−12 Watts per square meter) and an antenna has an effective area of 12 m2, then the antenna would deliver 12 pW of RF power to the receiver (30 microvolts RMS at 75 ohms). Since the receiving antenna is not equally sensitive to signals received from all directions, the effective area is a function of the direction to the source. Due to reciprocity (discussed above) the gain of an antenna used for transmitting must be proportional to its effective area when used for receiving. Consider an antenna with no loss, that is, one whose electrical efficiency is 100%. It can be shown that its effective area averaged over all directions must be equal to , the wavelength squared divided by . Gain is defined such that the average gain over all directions for an antenna with 100% electrical efficiency is equal to 1. Therefore, the effective area in terms of the gain in a given direction is given by: For an antenna with an efficiency of less than 100%, both the effective area and gain are reduced by that same amount. Therefore, the above relationship between gain and effective area still holds. These are thus two different ways of expressing the same quantity. eff is especially convenient when computing the power that would be received by an antenna of a specified gain, as illustrated by the above example. Radiation pattern The radiation pattern of an antenna is a plot of the relative field strength of the radio waves emitted by the antenna at different angles in the far field. It is typically represented by a three-dimensional graph, or polar plots of the horizontal and vertical cross sections. The pattern of an ideal isotropic antenna, which radiates equally in all directions, would look like a sphere. Many nondirectional antennas, such as monopoles and dipoles, emit equal power in all horizontal directions, with the power dropping off at higher and lower angles; this is called an omnidirectional pattern and when plotted looks like a torus or donut. The radiation of many antennas shows a pattern of maxima or "lobes" at various angles, separated by "nulls", angles where the radiation falls to zero. This is because the radio waves emitted by different parts of the antenna typically interfere, causing maxima at angles where the radio waves arrive at distant points in phase, and zero radiation at other angles where the radio waves arrive out of phase. In a directional antenna designed to project radio waves in a particular direction, the lobe in that direction is designed larger than the others and is called the "main lobe". The other lobes usually represent unwanted radiation and are called "sidelobes". The axis through the main lobe is called the "principal axis" or "boresight axis". The polar diagrams (and therefore the efficiency and gain) of Yagi antennas are tighter if the antenna is tuned for a narrower frequency range, e.g. the grouped antenna compared to the wideband. Similarly, the polar plots of horizontally polarized yagis are tighter than for those vertically polarized. Field regions The space surrounding an antenna can be divided into three concentric regions: The reactive near-field (also called the inductive near-field), the radiating near-field (Fresnel region) and the far-field (Fraunhofer) regions. These regions are useful to identify the field structure in each, although the transitions between them are gradual; there are no clear boundaries. The far-field region is far enough from the antenna to ignore its size and shape: It can be assumed that the electromagnetic wave is purely a radiating plane wave (electric and magnetic fields are in phase and perpendicular to each other and to the direction of propagation). This simplifies the mathematical analysis of the radiated field. Efficiency Efficiency of a transmitting antenna is the ratio of power actually radiated (in all directions) to the power absorbed by the antenna terminals. The power supplied to the antenna terminals which is not radiated is converted into heat. This is usually through loss resistance in the antenna's conductors, or loss between the reflector and feed horn of a parabolic antenna. Antenna efficiency is separate from impedance matching, which may also reduce the amount of power radiated using a given transmitter. If an SWR meter reads 150 W of incident power and 50 W of reflected power, that means 100 W have actually been absorbed by the antenna (ignoring transmission line losses). How much of that power has actually been radiated cannot be directly determined through electrical measurements at (or before) the antenna terminals, but would require (for instance) careful measurement of field strength. The loss resistance and efficiency of an antenna can be calculated once the field strength is known, by comparing it to the power supplied to the antenna. The loss resistance will generally affect the feedpoint impedance, adding to its resistive component. That resistance will consist of the sum of the radiation resistance rad and the loss resistance loss. If a current is delivered to the terminals of an antenna, then a power of 2 rad will be radiated and a power of 2 loss will be lost as heat. Therefore, the efficiency of an antenna is equal to . Only the total resistance rad + loss can be directly measured. According to reciprocity, the efficiency of an antenna used as a receiving antenna is identical to its efficiency as a transmitting antenna, described above. The power that an antenna will deliver to a receiver (with a proper impedance match) is reduced by the same amount. In some receiving applications, the very inefficient antennas may have little impact on performance. At low frequencies, for example, atmospheric or man-made noise can mask antenna inefficiency. For example, CCIR Rep. 258-3 indicates man-made noise in a residential setting at 40 MHz is about 28 dB above the thermal noise floor. Consequently, an antenna with a 20 dB loss (due to inefficiency) would have little impact on system noise performance. The loss within the antenna will affect the intended signal and the noise/interference identically, leading to no reduction in signal to noise ratio (SNR). Antennas which are not a significant fraction of a wavelength in size are inevitably inefficient due to their small radiation resistance. AM broadcast radios include a small loop antenna for reception which has an extremely poor efficiency. This has little effect on the receiver's performance, but simply requires greater amplification by the receiver's electronics. Contrast this tiny component to the massive and very tall towers used at AM broadcast stations for transmitting at the very same frequency, where every percentage point of reduced antenna efficiency entails a substantial cost. The definition of antenna gain or power gain already includes the effect of the antenna's efficiency. Therefore, if one is trying to radiate a signal toward a receiver using a transmitter of a given power, one need only compare the gain of various antennas rather than considering the efficiency as well. This is likewise true for a receiving antenna at very high (especially microwave) frequencies, where the point is to receive a signal which is strong compared to the receiver's noise temperature. However, in the case of a directional antenna used for receiving signals with the intention of rejecting interference from different directions, one is no longer concerned with the antenna efficiency, as discussed above. In this case, rather than quoting the antenna gain, one would be more concerned with the directive gain, or simply directivity which does not include the effect of antenna (in)efficiency. The directive gain of an antenna can be computed from the published gain divided by the antenna's efficiency. In equation form, gain = directivity × efficiency. Polarization The orientation and physical structure of an antenna determine the polarization of the electric field of the radio wave transmitted by it. For instance, an antenna composed of a linear conductor (such as a dipole or whip antenna) oriented vertically will result in vertical polarization; if turned on its side the same antenna's polarization will be horizontal. Reflections generally affect polarization. Radio waves reflected off the ionosphere can change the wave's polarization. For line-of-sight communications or ground wave propagation, horizontally or vertically polarized transmissions generally remain in about the same polarization state at the receiving location. Using a vertically polarized antenna to receive a horizontally polarized wave (or visa-versa) results in relatively poor reception. An antenna's polarization can sometimes be inferred directly from its geometry. When the antenna's conductors viewed from a reference location appear along one line, then the antenna's polarization will be linear in that very direction. In the more general case, the antenna's polarization must be determined through analysis. For instance, a turnstile antenna mounted horizontally (as is usual), from a distant location on Earth, appears as a horizontal line segment, so its radiation received there is horizontally polarized. But viewed at a downward angle from an airplane, the same antenna does not meet this requirement; in fact its radiation is elliptically polarized when viewed from that direction. In some antennas the state of polarization will change with the frequency of transmission. The polarization of a commercial antenna is an essential specification. In the most general case, polarization is elliptical, meaning that over each cycle the electric field vector traces out an ellipse. Two special cases are linear polarization (the ellipse collapses into a line) as discussed above, and circular polarization (in which the two axes of the ellipse are equal). In linear polarization the electric field of the radio wave oscillates along one direction. In circular polarization, the electric field of the radio wave rotates around the axis of propagation. Circular or elliptically polarized radio waves are designated as right-handed or left-handed using the "thumb in the direction of the propagation" rule. Note that for circular polarization, optical researchers use the opposite right-hand rule from the one used by radio engineers. It is best for the receiving antenna to match the polarization of the transmitted wave for optimum reception. Otherwise there will be a loss of signal strength: when a linearly polarized antenna receives linearly polarized radiation at a relative angle of θ, then there will be a power loss of cos2θ . A circularly polarized antenna can be used to equally well match vertical or horizontal linear polarizations, suffering a 3 dB signal reduction. However it will be blind to a circularly polarized signal of the opposite orientation. Impedance matching Maximum power transfer requires matching the impedance of an antenna system (as seen looking into the transmission line) to the complex conjugate of the impedance of the receiver or transmitter. In the case of a transmitter, however, the desired matching impedance might not exactly correspond to the dynamic output impedance of the transmitter as analyzed as a source impedance but rather the design value (typically 50 Ohms) required for efficient and safe operation of the transmitting circuitry. The intended impedance is normally resistive, but a transmitter (and some receivers) may have limited additional adjustments to cancel a certain amount of reactance, in order to "tweak" the match. When a transmission line is used in between the antenna and the transmitter (or receiver) one generally would like an antenna system whose impedance is resistive and nearly the same as the characteristic impedance of that transmission line, in addition to matching the impedance that the transmitter (or receiver) expects. The match is sought to minimize the amplitude of standing waves (measured via the standing wave ratio; SWR) that a mismatch raises on the line, and the increase in transmission line losses it entails. Antenna tuning at the antenna Antenna tuning, in the strict sense of modifying the antenna itself, generally refers only to cancellation of any reactance seen at the antenna terminals, leaving only a resistive impedance which might or might not be exactly the desired impedance (that of the available transmission line). Although an antenna may be designed to have a purely resistive feedpoint impedance (such as a dipole 97% of a half wavelength long) at just one frequency, this will very likely not be exactly true at other frequencies that the antenna is eventually used for. In most cases, in principle the physical length of the antenna can be "trimmed" to obtain a pure resistance, although this is rarely convenient. On the other hand, the addition of a contrary inductance or capacitance can be used to cancel a residual capacitive or inductive reactance, respectively, and may be more convenient than lowering and trimming or extending the antenna, then hoisting it back. Antenna reactance may be removed using lumped elements, such as capacitors or inductors in the main path of current traversing the antenna, often near the feedpoint, or by incorporating capacitive or inductive structures into the conducting body of the antenna to cancel the feedpoint reactance – such as open-ended "spoke" radial wires, or looped parallel wires – hence genuinely tune the antenna to resonance. In addition to those reactance-neutralizing add-ons, antennas of any kind may include a transformer and / or transformer balun at their feedpoint, to change the resistive part of the impedance to more nearly match the feedline's characteristic impedance. Line matching at the radio Antenna tuning in the loose sense, performed by an impedance matching device (somewhat inappropriately named an "antenna tuner", or the older, more appropriate term transmatch) goes beyond merely removing reactance and includes transforming the remaining resistance to match the feedline and radio. An additional problem is matching the remaining resistive impedance to the characteristic impedance of the transmission line: A general impedance matching network (an "antenna tuner" or ATU) will have at least two adjustable elements to correct both components of impedance. Any matching network will have both power losses and power restrictions when used for transmitting. Commercial antennas are generally designed to approximately match standard 50 Ohm coaxial cables, at standard frequencies; the design expectation is that a matching network will be merely used to 'tweak' any residual mismatch. Extreme examples of loaded small antennas In some cases matching is done in a more extreme manner, not simply to cancel a small amount of residual reactance, but to resonate an antenna whose resonance frequency is quite different from the intended frequency of operation. Short vertical "whip" For instance, for practical reasons a "whip antenna" can be made significantly shorter than a quarter-wavelength and then resonated, using a so-called loading coil. The physically large inductor at the base of the antenna has an inductive reactance which is the opposite of the capacitative reactance that the short vertical antenna has at the desired operating frequency. The result is a pure resistance seen at feedpoint of the loading coil; although, without further measures, the resistance will be somewhat lower than would be desired to match commercial coax. Small "magnetic" loop Another extreme case of impedance matching occurs when using a small loop antenna (usually, but not always, for receiving) at a relatively low frequency, where it appears almost as a pure inductor. When such an inductor is resonated via a capacitor attached in parallel across its feedpoint, the capacitor not only cancels the reactance but also greatly magnifies the very small radiation resistance of a small loop to produce a better-matched feedpoint resistance. This is the type of antenna used in most portable AM broadcast receivers (other than car radios): The standard AM antenna is a loop of wire wound around a ferrite rod (a "loopstick antenna"). The loop is resonated by a coupled tuning capacitor, which is configured to match the receiver's tuning, in order to keep the antenna resonant at the chosen receive frequency over the AM broadcast band. Effect of ground Ground reflections is one of the common types of multipath. The radiation pattern and even the driving point impedance of an antenna can be influenced by the dielectric constant and especially conductivity of nearby objects. For a terrestrial antenna, the ground is usually one such object of importance. The antenna's height above the ground, as well as the electrical properties (permittivity and conductivity) of the ground, can then be important. Also, in the particular case of a monopole antenna, the ground (or an artificial ground plane) serves as the return connection for the antenna current thus having an additional effect, particularly on the impedance seen by the feed line. When an electromagnetic wave strikes a plane surface such as the ground, part of the wave is transmitted into the ground and part of it is reflected, according to the Fresnel coefficients. If the ground is a very good conductor then almost all of the wave is reflected (180° out of phase), whereas a ground modeled as a (lossy) dielectric can absorb a large amount of the wave's power. The power remaining in the reflected wave, and the phase shift upon reflection, strongly depend on the wave's angle of incidence and polarization. The dielectric constant and conductivity (or simply the complex dielectric constant) is dependent on the soil type and is a function of frequency. For very low frequencies to high frequencies (< 30 MHz), the ground behaves as a lossy dielectric, thus the ground is characterized both by a conductivity and permittivity (dielectric constant) which can be measured for a given soil (but is influenced by fluctuating moisture levels) or can be estimated from certain maps. At lower mediumwave frequencies the ground acts mainly as a good conductor, which AM broadcast (0.5–1.7 MHz) antennas depend on. At frequencies between 3–30 MHz, a large portion of the energy from a horizontally polarized antenna reflects off the ground, with almost total reflection at the grazing angles important for ground wave propagation. That reflected wave, with its phase reversed, can either cancel or reinforce the direct wave, depending on the antenna height in wavelengths and elevation angle (for a sky wave). On the other hand, vertically polarized radiation is not well reflected by the ground except at grazing incidence or over very highly conducting surfaces such as sea water. However the grazing angle reflection important for ground wave propagation, using vertical polarization, is in phase with the direct wave, providing a boost of up to 6 dB, as is detailed below. At VHF and above (> 30 MHz) the ground becomes a poorer reflector. However, for shortwave frequencies, especially below ~15 MHz, it remains a good reflector especially for horizontal polarization and grazing angles of incidence. That is important as these higher frequencies usually depend on horizontal line-of-sight propagation (except for satellite communications), the ground then behaving almost as a mirror. The net quality of a ground reflection depends on the topography of the surface. When the irregularities of the surface are much smaller than the wavelength, the dominant regime is that of specular reflection, and the receiver sees both the real antenna and an image of the antenna under the ground due to reflection. But if the ground has irregularities not small compared to the wavelength, reflections will not be coherent but shifted by random phases. With shorter wavelengths (higher frequencies), this is generally the case. Whenever both the receiving or transmitting antenna are placed at significant heights above the ground (relative to the wavelength), waves reflected specularly by the ground will travel a longer distance than direct waves, inducing a phase shift which can sometimes be significant. When a sky wave is launched by such an antenna, that phase shift is always significant unless the antenna is very close to the ground (compared to the wavelength). The phase of reflection of electromagnetic waves depends on the polarization of the incident wave. Given the larger refractive index of the ground (typically ≈ 2) compared to air ( = 1), the phase of horizontally polarized radiation is reversed upon reflection (a phase shift of radians, or 180°). On the other hand, the vertical component of the wave's electric field is reflected at grazing angles of incidence approximately in phase. These phase shifts apply as well to a ground modeled as a good electrical conductor. This means that a receiving antenna "sees" an image of the emitting antenna but with 'reversed' currents (opposite in direction and phase) if the emitting antenna is horizontally oriented (and thus horizontally polarized). However, the received current will be in the same absolute direction and phase if the emitting antenna is vertically polarized. The actual antenna which is transmitting the original wave then also may receive a strong signal from its own image from the ground. This will induce an additional current in the antenna element, changing the current at the feedpoint for a given feedpoint voltage. Thus the antenna's impedance, given by the ratio of feedpoint voltage to current, is altered due to the antenna's proximity to the ground. This can be quite a significant effect when the antenna is within a wavelength or two of the ground. But as the antenna height is increased, the reduced power of the reflected wave (due to the inverse square law) allows the antenna to approach its asymptotic feedpoint impedance given by theory. At lower heights, the effect on the antenna's impedance is very sensitive to the exact distance from the ground, as this affects the phase of the reflected wave relative to the currents in the antenna. Changing the antenna's height by a quarter wavelength, then changes the phase of the reflection by 180°, with a completely different effect on the antenna's impedance. The ground reflection has an important effect on the net far field radiation pattern in the vertical plane, that is, as a function of elevation angle, which is thus different between a vertically and horizontally polarized antenna. Consider an antenna at a height above the ground, transmitting a wave considered at the elevation angle . For a vertically polarized transmission the magnitude of the electric field of the electromagnetic wave produced by the direct ray plus the reflected ray is: Thus the power received can be as high as 4 times that due to the direct wave alone (such as when = 0), following the square of the cosine. The sign inversion for the reflection of horizontally polarized emission instead results in: where: is the electrical field that would be received by the direct wave if there were no ground. is the elevation angle of the wave being considered. is the wavelength. is the height of the antenna (half the distance between the antenna and its image). For horizontal propagation between transmitting and receiving antennas situated near the ground reasonably far from each other, the distances traveled by the direct and reflected rays are nearly the same. There is almost no relative phase shift. If the emission is polarized vertically, the two fields (direct and reflected) add and there is maximum of received signal. If the signal is polarized horizontally, the two signals subtract and the received signal is largely cancelled. The vertical plane radiation patterns are shown in the image at right. With vertical polarization there is always a maximum for = 0, horizontal propagation (left pattern). For horizontal polarization, there is cancellation at that angle. The above formulae and these plots assume the ground as a perfect conductor. These plots of the radiation pattern correspond to a distance between the antenna and its image of 2.5 . As the antenna height is increased, the number of lobes increases as well. The difference in the above factors for the case of = 0 is the reason that most broadcasting (transmissions intended for the public) uses vertical polarization. For receivers near the ground, horizontally polarized transmissions suffer cancellation. For best reception the receiving antennas for these signals are likewise vertically polarized. In some applications where the receiving antenna must work in any position, as in mobile phones, the base station antennas use mixed polarization, such as linear polarization at an angle (with both vertical and horizontal components) or circular polarization. On the other hand, analog television transmissions are usually horizontally polarized, because in urban areas buildings can reflect the electromagnetic waves and create ghost images due to multipath propagation. Using horizontal polarization, ghosting is reduced because the amount of reflection in the horizontal polarization off the side of a building is generally less than in the vertical direction. Vertically polarized analog television have been used in some rural areas. In digital terrestrial television such reflections are less problematic, due to robustness of binary transmissions and error correction. Modeling antennas with line equations In the first approximation, the current in a thin antenna is distributedexactly as in a transmission line. — Schelkunoff & Friis (1952) The flow of current in wire antennas is identical to the solution of counter-propagating waves in a single conductor transmission line, which can be solved using the telegrapher's equations. Solutions of currents along antenna elements are more conveniently and accurately obtained by numerical methods, so transmission-line techniques have largely been abandoned for precision modelling, but they continue to be a widely used source of useful, simple approximations that describe well the impedance profiles of antennas. Unlike transmission lines, currents in antennas contribute power to the radiated part electromagnetic field, which can be modeled using radiation resistance. The end of an antenna element corresponds to an unterminated (open) end of a single-conductor transmission line, resulting in a reflected wave identical to the incident wave, with its voltage in phase with the incident wave and its current in the opposite phase (thus net zero current, where there is, after all, no further conductor). The combination of the incident and reflected wave, just as in a transmission line, forms a standing wave with a current node at the conductor's end, and a voltage node one-quarter wavelength from the end (if the element is at least that long). In a resonant antenna, the feedpoint of the antenna is at one of those voltage nodes. Due to discrepancies from the simplified version of the transmission line model, the voltage one quarter wavelength from the current node is not exactly zero, but it is near a minimum, and small compared to the much large voltage at the conductor's end. Hence, a feed point matching the antenna at that spot requires a relatively small voltage but large current (the currents from the two waves add in-phase there), thus a relatively low feedpoint impedance. Feeding the antenna at other points involves a large voltage, thus a large impedance, and usually one that is primarily reactive (low power factor), which is a terrible impedance match to available transmission lines. Therefore, it is usually desired for an antenna to operate as a resonant element with each conductor having a length of one quarter wavelength (or any other odd multiples of a quarter wavelength). For instance, a half-wave dipole has two such elements (one connected to each conductor of a balanced transmission line) about one quarter wavelength long. Depending on the conductors' diameters, a small deviation from this length is adopted in order to reach the point where the antenna current and the (small) feedpoint voltage are exactly in phase. Then the antenna presents a purely resistive impedance, and ideally one close to the characteristic impedance of an available transmission line. Despite these useful properties, resonant antennas have the disadvantage that they achieve resonance (purely resistive feedpoint impedance) only at a fundamental frequency, and perhaps some of its harmonics, and the feedpoint resistance is larger at higher-order resonances. Therefore, resonant antennas can only achieve their good performance within a limited bandwidth, depending on the at the resonance. Mutual impedance and interaction between antennas The electric and magnetic fields emanating from a driven antenna element will generally affect the voltages and currents in nearby antennas, antenna elements, or other conductors. This is particularly true when the affected conductor is a resonant element (multiple of half-wavelengths in length) at about the same frequency, as is the case where the conductors are all part of the same active or passive antenna array. Because the affected conductors are in the near-field, one can not just treat two antennas as transmitting and receiving a signal according to the Friis transmission formula for instance, but must calculate the mutual impedance matrix which takes into account both voltages and currents (interactions through both the electric and magnetic fields). Thus using the mutual impedances calculated for a specific geometry, one can solve for the radiation pattern of a Yagi–Uda antenna or the currents and voltages for each element of a phased array. Such an analysis can also describe in detail reflection of radio waves by a ground plane or by a corner reflector and their effect on the impedance (and radiation pattern) of an antenna in its vicinity. Often such near-field interactions are undesired and pernicious. Currents in random metal objects near a transmitting antenna will often be in poor conductors, causing loss of RF power in addition to unpredictably altering the characteristics of the antenna. By careful design, it is possible to reduce the electrical interaction between nearby conductors. For instance, the 90 degree angle in between the two dipoles composing the turnstile antenna insures no interaction between these, allowing these to be driven independently (but actually with the same signal in quadrature phases in the turnstile antenna design). Antenna types Antennas can be classified by operating principles or by their application. Different authorities placed antennas in narrower or broader categories. Generally these include Dipole and monopole antennas Array antennas Loop antennas Aperture antennas Traveling wave antennas Log-periodic antenna Spiral antenna Horn antenna Adcock antenna Sector antenna Helical antenna These antenna types and others are summarized in greater detail in the overview article, Antenna types, as well as in each of the linked articles in the list above, and in even more detail in articles which those link to.	Technology	Components	null
187344	https://en.wikipedia.org/wiki/Oil%20drop%20experiment	Oil drop experiment	The oil drop experiment was performed by Robert A. Millikan and Harvey Fletcher in 1909 to measure the elementary electric charge (the charge of the electron). The experiment took place in the Ryerson Physical Laboratory at the University of Chicago. Millikan received the Nobel Prize in Physics in 1923. The experiment observed tiny electrically charged droplets of oil located between two parallel metal surfaces, forming the plates of a capacitor. The plates were oriented horizontally, with one plate above the other. A mist of atomized oil drops was introduced through a small hole in the top plate and was ionized by x-rays, making them negatively charged. First, with zero applied electric field, the velocity of a falling droplet was measured. At terminal velocity, the drag force equals the gravitational force. As both forces depend on the radius in different ways, the radius of the droplet, and therefore the mass and gravitational force, could be determined (using the known density of the oil). Next, a voltage inducing an electric field was applied between the plates and adjusted until the drops were suspended in mechanical equilibrium, indicating that the electrical force and the gravitational force were in balance. Using the known electric field, Millikan and Fletcher could determine the charge on the oil droplet. By repeating the experiment for many droplets, they confirmed that the charges were all small integer multiples of a certain base value, which was found to be , about 0.6% difference from the currently accepted value of They proposed that this was the magnitude of the negative charge of a single electron. Background Starting in 1908, while a professor at the University of Chicago, Millikan, with the significant input of Fletcher, the "able assistance of Mr. J. Yinbong Lee", and after improving his setup, published his seminal study in 1913. This remains controversial since papers found after Fletcher's death describe events in which Millikan coerced Fletcher into relinquishing authorship as a condition for receiving his PhD. In return, Millikan used his influence in support of Fletcher's career at Bell Labs. Millikan and Fletcher's experiment involved measuring the force on oil droplets in a glass chamber sandwiched between two electrodes, one above and one below. With the electrical field calculated, they could measure the droplet's charge, the charge on a single electron being (). At the time of Millikan and Fletcher's oil drop experiments, the existence of subatomic particles was not universally accepted. Experimenting with cathode rays in 1897, J. J. Thomson had discovered negatively charged "corpuscles", as he called them, with a mass about 1/1837 that of a hydrogen atom. Similar results had been found by George FitzGerald and Walter Kaufmann. Most of what was then known about electricity and magnetism, however, could be explained on the basis that charge is a continuous variable; in much the same way that many of the properties of light can be explained by treating it as a continuous wave rather than as a stream of photons. The elementary charge e is one of the fundamental physical constants and thus the accuracy of the value is of great importance. In 1923, Millikan won the Nobel Prize in physics, in part because of this experiment. Thomas Edison, who had previously thought of charge as a continuous variable, became convinced after working with Millikan and Fletcher's apparatus. This experiment has since been repeated by generations of physics students, although it is rather expensive and difficult to conduct properly. From 1995 to 2007, several computer-automated experiments have been conducted at SLAC to search for isolated fractionally charged particles, however, no evidence for fractional charge particles has been found after measuring over 100 million drops. Experimental procedure Apparatus Millikan's and Fletcher's apparatus incorporated a parallel pair of horizontal metal plates. By applying a potential difference across the plates, a uniform electric field was created in the space between them. A ring of insulating material was used to hold the plates apart. Four holes were cut into the ring, three for illumination by a bright light, and another to allow viewing through a microscope. A fine mist of oil droplets was sprayed into a chamber above the plates. The oil was of a type usually used in vacuum apparatus and was chosen because it had an extremely low vapour pressure. Ordinary oils would evaporate under the heat of the light source causing the mass of the oil drop to change over the course of the experiment. Some oil drops became electrically charged through friction with the nozzle as they were sprayed. Alternatively, charging could be brought about by including an ionizing radiation source (such as an X-ray tube). The droplets entered the space between the plates and, because they were charged, could be made to rise and fall by changing the voltage across the plates. Method Initially the oil drops are allowed to fall between the plates with the electric field turned off. They very quickly reach a terminal velocity because of friction with the air in the chamber. The field is then turned on and, if it is large enough, some of the drops (the charged ones) will start to rise. (This is because the upwards electric force FE is greater for them than the downwards gravitational force Fg, in the same way bits of paper can be picked by a charged rubber rod). A likely looking drop is selected and kept in the middle of the field of view by alternately switching off the voltage until all the other drops have fallen. The experiment is then continued with this one drop. The drop is allowed to fall and its terminal velocity v1 in the absence of an electric field is calculated. The drag force acting on the drop can then be worked out using Stokes' law: where v1 is the terminal velocity (i.e. velocity in the absence of an electric field) of the falling drop, η is the viscosity of the air, and r is the radius of the drop. The weight w is the volume D multiplied by the density ρ and the acceleration due to gravity g. However, what is needed is the apparent weight. The apparent weight in air is the true weight minus the upthrust (which equals the weight of air displaced by the oil drop). For a perfectly spherical droplet the apparent weight can be written as: At terminal velocity the oil drop is not accelerating. Therefore, the total force acting on it must be zero and the two forces F and must cancel one another out (that is, ). This implies Once r is calculated, can easily be worked out. Now the field is turned back on, and the electric force on the drop is where q is the charge on the oil drop and E is the electric field between the plates. For parallel plates where V is the potential difference and d is the distance between the plates. One conceivable way to work out q would be to adjust V until the oil drop remained steady. Then we could equate FE with . Also, determining FE proves difficult because the mass of the oil drop is difficult to determine without reverting to the use of Stokes' Law. A more practical approach is to turn V up slightly so that the oil drop rises with a new terminal velocity v2. Then Comparison to modern values Effective from the 2019 revision of the SI, the value of the elementary charge is defined to be exactly . Before that, the most recent (2014) accepted value was , where the (98) indicates the uncertainty of the last two decimal places. In his Nobel lecture, Millikan gave his measurement as , which equals . The difference is less than one percent, but is six times greater than Millikan's standard error, so the disagreement is significant. Using X-ray experiments, Erik Bäcklin in 1928 found a higher value of the elementary charge, or , which is within uncertainty of the exact value. Raymond Thayer Birge, conducting a review of physical constants in 1929, stated "The investigation by Bäcklin constitutes a pioneer piece of work, and it is quite likely, as such, to contain various unsuspected sources of systematic error. If [... it is ...] weighted according to the apparent probable error [...], the weighted average will still be suspiciously high. [...] the writer has finally decided to reject the Bäcklin value, and to use the weighted mean of the remaining two values." Birge averaged Millikan's result and a different, less accurate X-ray experiment that agreed with Millikan's result. Successive X-ray experiments continued to give high results, and proposals for the discrepancy were ruled out experimentally. Sten von Friesen measured the value with a new electron diffraction method, and the oil drop experiment was redone. Both gave high numbers. By 1937 it was "quite obvious" that Millikan's value could not be maintained any longer, and the established value became or . Controversy Some controversy was raised by physicist Gerald Holton (1978) who pointed out that Millikan recorded more measurements in his journal than he included in his final results. Holton suggested these data points were omitted from the large set of oil drops measured in his experiments without apparent reason. This claim was disputed by Allan Franklin, a high energy physics experimentalist and philosopher of science at the University of Colorado. Franklin contended that Millikan's exclusions of data did not substantively affect his final value of e, but did reduce the statistical error around this estimate e. This enabled Millikan to claim that he had calculated e to better than one half of one percent; in fact, if Millikan had included all of the data he had thrown out, the standard error of the mean would have been within 2%. While this would still have resulted in Millikan having measured e better than anyone else at the time, the slightly larger uncertainty might have allowed more disagreement with his results within the physics community. While Franklin left his support for Millikan's measurement with the conclusion that concedes that Millikan may have performed "cosmetic surgery" on the data, David Goodstein investigated the original detailed notebooks kept by Millikan, concluding that Millikan plainly states here and in the reports that he included only drops that had undergone a "complete series of observations" and excluded no drops from this group of complete measurements. Reasons for a failure to generate a complete observation include annotations regarding the apparatus setup, oil drop production, and atmospheric effects which invalidated, in Millikan's opinion (borne out by the reduced error in this set), a given particular measurement. Millikan's experiment as an example of psychological effects in scientific methodology In a commencement address given at the California Institute of Technology (Caltech) in 1974 (and reprinted in Surely You're Joking, Mr. Feynman! in 1985 as well as in The Pleasure of Finding Things Out in 1999), physicist Richard Feynman noted:	Physical sciences	Basics_9	null
187360	https://en.wikipedia.org/wiki/Magnetic%20susceptibility	Magnetic susceptibility	In electromagnetism, the magnetic susceptibility (; denoted , chi) is a measure of how much a material will become magnetized in an applied magnetic field. It is the ratio of magnetization (magnetic moment per unit volume) to the applied magnetic field intensity . This allows a simple classification, into two categories, of most materials' responses to an applied magnetic field: an alignment with the magnetic field, , called paramagnetism, or an alignment against the field, , called diamagnetism. Magnetic susceptibility indicates whether a material is attracted into or repelled out of a magnetic field. Paramagnetic materials align with the applied field and are attracted to regions of greater magnetic field. Diamagnetic materials are anti-aligned and are pushed away, toward regions of lower magnetic fields. On top of the applied field, the magnetization of the material adds its own magnetic field, causing the field lines to concentrate in paramagnetism, or be excluded in diamagnetism. Quantitative measures of the magnetic susceptibility also provide insights into the structure of materials, providing insight into bonding and energy levels. Furthermore, it is widely used in geology for paleomagnetic studies and structural geology. The magnetizability of materials comes from the atomic-level magnetic properties of the particles of which they are made. Usually, this is dominated by the magnetic moments of electrons. Electrons are present in all materials, but without any external magnetic field, the magnetic moments of the electrons are usually either paired up or random so that the overall magnetism is zero (the exception to this usual case is ferromagnetism). The fundamental reasons why the magnetic moments of the electrons line up or do not are very complex and cannot be explained by classical physics. However, a useful simplification is to measure the magnetic susceptibility of a material and apply the macroscopic form of Maxwell's equations. This allows classical physics to make useful predictions while avoiding the underlying quantum mechanical details. Definition Volume susceptibility Magnetic susceptibility is a dimensionless proportionality constant that indicates the degree of magnetization of a material in response to an applied magnetic field. A related term is magnetizability, the proportion between magnetic moment and magnetic flux density. A closely related parameter is the permeability, which expresses the total magnetization of material and volume. The volume magnetic susceptibility, represented by the symbol (often simply , sometimes – magnetic, to distinguish from the electric susceptibility), is defined in the International System of Units – in other systems there may be additional constants – by the following relationship: Here, is the magnetization of the material (the magnetic dipole moment per unit volume), with unit amperes per meter, and is the magnetic field strength, also with the unit amperes per meter. is therefore a dimensionless quantity. Using SI units, the magnetic induction is related to by the relationship where is the vacuum permeability (see table of physical constants), and is the relative permeability of the material. Thus the volume magnetic susceptibility and the magnetic permeability are related by the following formula: Sometimes an auxiliary quantity called intensity of magnetization (also referred to as magnetic polarisation ) and with unit teslas, is defined as This allows an alternative description of all magnetization phenomena in terms of the quantities and , as opposed to the commonly used and . Molar susceptibility and mass susceptibility There are two other measures of susceptibility, the molar magnetic susceptibility () with unit m3/mol, and the mass magnetic susceptibility () with unit m3/kg that are defined below, where is the density with unit kg/m3 and is molar mass with unit kg/mol: In CGS units The definitions above are according to the International System of Quantities (ISQ) upon which the SI is based. However, many tables of magnetic susceptibility give the values of the corresponding quantities of the CGS system (more specifically CGS-EMU, short for electromagnetic units, or Gaussian-CGS; both are the same in this context). The quantities characterizing the permeability of free space for each system have different defining equations: The respective CGS susceptibilities are multiplied by 4 to give the corresponding ISQ quantities (often referred to as SI quantities) with the same units: For example, the CGS volume magnetic susceptibility of water at 20 °C is , which is using the SI convention, both quantities being dimensionless. Whereas for most electromagnetic quantities, which system of quantities it belongs to can be disambiguated by incompatibility of their units, this is not true for the susceptibility quantities. In physics it is common to see CGS mass susceptibility with unit cm3/g or emu/g⋅Oe−1, and the CGS molar susceptibility with unit cm3/mol or emu/mol⋅Oe−1. Paramagnetism and diamagnetism If is positive, a material can be paramagnetic. In this case, the magnetic field in the material is strengthened by the induced magnetization. Alternatively, if is negative, the material is diamagnetic. In this case, the magnetic field in the material is weakened by the induced magnetization. Generally, nonmagnetic materials are said to be para- or diamagnetic because they do not possess permanent magnetization without external magnetic field. Ferromagnetic, ferrimagnetic, or antiferromagnetic materials possess permanent magnetization even without external magnetic field and do not have a well defined zero-field susceptibility. Experimental measurement Volume magnetic susceptibility is measured by the force change felt upon a substance when a magnetic field gradient is applied. Early measurements are made using the Gouy balance where a sample is hung between the poles of an electromagnet. The change in weight when the electromagnet is turned on is proportional to the susceptibility. Today, high-end measurement systems use a superconductive magnet. An alternative is to measure the force change on a strong compact magnet upon insertion of the sample. This system, widely used today, is called the Evans balance. For liquid samples, the susceptibility can be measured from the dependence of the NMR frequency of the sample on its shape or orientation. Another method using NMR techniques measures the magnetic field distortion around a sample immersed in water inside an MR scanner. This method is highly accurate for diamagnetic materials with susceptibilities similar to water. Tensor susceptibility The magnetic susceptibility of most crystals is not a scalar quantity. Magnetic response is dependent upon the orientation of the sample and can occur in directions other than that of the applied field . In these cases, volume susceptibility is defined as a tensor: where and refer to the directions (e.g., of the and Cartesian coordinates) of the applied field and magnetization, respectively. The tensor is thus degree 2 (second order), dimension (3,3) describing the component of magnetization in the th direction from the external field applied in the th direction. Differential susceptibility In ferromagnetic crystals, the relationship between and is not linear. To accommodate this, a more general definition of differential susceptibility is used: where is a tensor derived from partial derivatives of components of with respect to components of . When the coercivity of the material parallel to an applied field is the smaller of the two, the differential susceptibility is a function of the applied field and self interactions, such as the magnetic anisotropy. When the material is not saturated, the effect will be nonlinear and dependent upon the domain wall configuration of the material. Several experimental techniques allow for the measurement of the electronic properties of a material. An important effect in metals under strong magnetic fields, is the oscillation of the differential susceptibility as function of . This behaviour is known as the De Haas–Van Alphen effect and relates the period of the susceptibility with the Fermi surface of the material. An analogue non-linear relation between magnetization and magnetic field happens for antiferromagnetic materials. In the frequency domain When the magnetic susceptibility is measured in response to an AC magnetic field (i.e. a magnetic field that varies sinusoidally), this is called AC susceptibility. AC susceptibility (and the closely related "AC permeability") are complex number quantities, and various phenomena, such as resonance, can be seen in AC susceptibility that cannot occur in constant-field (DC) susceptibility. In particular, when an AC field is applied perpendicular to the detection direction (called the "transverse susceptibility" regardless of the frequency), the effect has a peak at the ferromagnetic resonance frequency of the material with a given static applied field. Currently, this effect is called the microwave permeability or network ferromagnetic resonance in the literature. These results are sensitive to the domain wall configuration of the material and eddy currents. In terms of ferromagnetic resonance, the effect of an AC-field applied along the direction of the magnetization is called parallel pumping. Table of examples Sources of published data The CRC Handbook of Chemistry and Physics has one of the few published magnetic susceptibility tables. The data are listed as CGS quantities. The molar susceptibility of several elements and compounds are listed in the CRC. Application in the geosciences In Earth science, magnetism is a useful parameter to describe and analyze rocks. Additionally, the anisotropy of magnetic susceptibility (AMS) within a sample determines parameters as directions of paleocurrents, maturity of paleosol, flow direction of magma injection, tectonic strain, etc. It is a non-destructive tool which quantifies the average alignment and orientation of magnetic particles within a sample.	Physical sciences	Magnetostatics	Physics
187377	https://en.wikipedia.org/wiki/Shipbuilding	Shipbuilding	Shipbuilding is the construction of ships and other floating vessels. In modern times, it normally takes place in a specialized facility known as a shipyard. Shipbuilders, also called shipwrights, follow a specialized occupation that traces its roots to before recorded history. Until recently, with the development of complex non-maritime technologies, a ship has often represented the most advanced structure that the society building it could produce. Some key industrial advances were developed to support shipbuilding, for instance the sawing of timbers by mechanical saws propelled by windmills in Dutch shipyards during the first half of the 17th century. The design process saw the early adoption of the logarithm (invented in 1615) to generate the curves used to produce the shape of a hull, especially when scaling up these curves accurately in the mould loft. Shipbuilding and ship repairs, both commercial and military, are referred to as naval engineering. The construction of boats is a similar activity called boat building. The dismantling of ships is called ship breaking. The earliest evidence of maritime transport by modern humans is the settlement of Australia between 50,000 and 60,000 years ago. This almost certainly involved rafts, possibly equipped with some sort of sail. Much of the development beyond that raft technology occurred in the "nursery" areas of the Mediterranean and in Maritime Southeast Asia. Favoured by warmer waters and a number of inter-visible islands, boats (and, later, ships) with water-tight hulls (unlike the "flow through" structure of a raft) could be developed. The ships of ancient Egypt were built by joining the hull planks together, edge to edge, with tenons set in mortices cut in the mating edges. A similar technique, but with the tenons being pinned in position by dowels, was used in the Mediterranean for most of classical antiquity. Both these variants are "shell first" techniques, where any reinforcing frames are inserted after assembly of the planking has defined the hull shape. Carvel construction then took over in the Mediterranean. Northern Europe used clinker construction, but with some flush-planked ship-building in, for instance, the bottom planking of cogs. The north-European and Mediterranean traditions merged in the late 15th century, with carvel construction being adopted in the North and the centre-line mounted rudder replacing the quarter rudder of the Mediterranean. These changes broadly coincided with improvements in sailing rigs, with the three masted ship becoming common, with square sails on the fore and main masts, and a fore and aft sail on the mizzen. Ship-building then saw a steady improvement in design techniques and introduction of new materials. Iron was used for more than fastenings (nails and bolts) as structural components such as iron knees were introduced, with examples existing in the mid-18th century and from the mid-19th century onwards. This was partly led by the shortage of "compass timber", the naturally curved timber that meant that shapes could be cut without weaknesses caused by cuts across the grain of the timber. Ultimately, whole ships were made of iron and, later, steel. History Pre-history The earliest known depictions (including paintings and models) of shallow-water sailing boats is from the 6th to 5th millennium BC of the Ubaid period of Mesopotamia. They were made from bundled reeds coated in bitumen and had bipod masts. They sailed in shallow coastal waters of the Persian Gulf. 4th millennium BC Ancient Egypt Evidence from Ancient Egypt shows that the early Egyptians knew how to assemble planks of wood into a ship hull as early as 3100 BC. Egyptian pottery as old as 4000 BC shows designs of early fluvial boats or other means for navigation. The Archaeological Institute of America reports that some of the oldest ships yet unearthed are known as the Abydos boats. These are a group of 14 ships discovered in Abydos that were constructed of wooden planks which were "sewn" together. Discovered by Egyptologist David O'Connor of New York University, woven straps were found to have been used to lash the planks together, and reeds or grass stuffed between the planks helped to seal the seams. Because the ships are all buried together and near a mortuary belonging to Pharaoh Khasekhemwy, originally they were all thought to have belonged to him, but one of the 14 ships dates to 3000 BC, and the associated pottery jars buried with the vessels also suggest earlier dating. The ship dating to 3000 BC was about 75 feet (23 m) long and is now thought to perhaps have belonged to an earlier pharaoh. According to professor O'Connor, the 5,000-year-old ship may have even belonged to Pharaoh Aha. Austronesia The Austronesian expansion, which began with migration from Taiwan to the island of Luzon in the Philippines, spread across Island Southeast Asia. Then, between 1500 BC and 1500 AD they settled uninhabited islands of the Pacific, and also sailed westward to Madagascar. This is associated with distinctive maritime technology: lashed lug construction techniques (both in outrigger canoes and in large planked sailing vessels), various types of outrigger and twin-hulled canoes and a range of sailing rigs that included the crab claw sail. The origins of this technology is difficult to date, relying largely on linguistics (studying the words for parts of boats), the written comments of people from other cultures, including the observations of European explorers at the time of first contact and the later more systematic ethnographic observations of the types of craft in use. There is only a small body of archaeological evidence available. Since Island Southeast Asia contained effective maritime transport between its very large number of islands long before Austronesian seafaring, it is argued that Austronesians adopted an existing maritime technology from the existing inhabitants of this region. Austronesian ships varied from simple canoes to large multihull ships. The simplest form of all ancestral Austronesian boats had five parts. The bottom part consists of a single piece of hollowed-out log. At the sides were two planks, and two horseshoe-shaped wood pieces formed the prow and stern. These were fitted tightly together edge-to-edge with dowels inserted into holes in between, and then lashed to each other with ropes (made from rattan or fiber) wrapped around protruding lugs on the planks. This characteristic and ancient Austronesian boatbuilding practice is known as the "lashed-lug" technique. They were commonly caulked with pastes made from various plants as well as tapa bark and fibres which would expand when wet, further tightening joints and making the hull watertight. They formed the shell of the boat, which was then reinforced by horizontal ribs. Shipwrecks of Austronesian ships can be identified from this construction as well as the absence of metal nails. Austronesian ships traditionally had no central rudders but were instead steered using an oar on one side. Austronesians traditionally made their sails from woven mats of the resilient and salt-resistant pandanus leaves. These sails allowed Austronesians to embark on long-distance voyaging. The ancient Champa of Vietnam also uniquely developed basket-hulled boats whose hulls were composed of woven and resin-caulked bamboo, either entirely or in conjunction with plank strakes. They range from small coracles (the o thúng) to large ocean-going trading ships like the ghe mành. 3rd millennium BC Ancient Egypt Early Egyptians also knew how to assemble planks of wood with treenails to fasten them together, using pitch for caulking the seams. The "Khufu ship", a 43.6-meter vessel sealed into a pit in the Giza pyramid complex at the foot of the Great Pyramid of Giza in the Fourth Dynasty around 2500 BC, is a full-size surviving example which may have fulfilled the symbolic function of a solar barque. Early Egyptians also knew how to fasten the planks of this ship together with mortise and tenon joints. Indus Valley The oldest known tidal dock in the world was built around 2500 BC during the Harappan civilisation at Lothal near the present day Mangrol harbour on the Gujarat coast in India. Other ports were probably at Balakot and Dwarka. However, it is probable that many small-scale ports, and not massive ports, were used for the Harappan maritime trade. Ships from the harbour at these ancient port cities established trade with Mesopotamia. Shipbuilding and boatmaking may have been prosperous industries in ancient India. Native labourers may have manufactured the flotilla of boats used by Alexander the Great to navigate across the Hydaspes and even the Indus, under Nearchos. The Indians also exported teak for shipbuilding to ancient Persia. Other references to Indian timber used for shipbuilding is noted in the works of Ibn Jubayr. 2nd millennium BC Mediterranean The ships of Ancient Egypt's Eighteenth Dynasty were typically about 25 meters (80 ft) in length and had a single mast, sometimes consisting of two poles lashed together at the top making an "A" shape. They mounted a single square sail on a yard, with an additional spar along the bottom of the sail. These ships could also be oar propelled. The ocean- and sea-going ships of Ancient Egypt were constructed with cedar wood, most likely hailing from Lebanon. The ships of Phoenicia seem to have been of a similar design. 1st millennium BC Austronesia Austronesians established the Austronesian maritime trade network at around 1000 to 600 BC, linking Southeast Asia with East Asia, South Asia, the Middle East, and later East Africa. The route later became part of the Spice trade network and the Maritime Silk Road. China The naval history of China stems back to the Spring and Autumn period (722 BC–481 BC) of the ancient Chinese Zhou dynasty. The Chinese built large rectangular barges known as "castle ships", which were essentially floating fortresses complete with multiple decks with guarded ramparts. However, the Chinese vessels during this era were essentially fluvial (riverine). True ocean-going Chinese fleets did not appear until the 10th century Song dynasty. Mediterranean There is considerable knowledge regarding shipbuilding and seafaring in the ancient Mediterranean. 1st millennium AD Austronesia Large multi-masted seafaring ships of Southeast Asian Austronesians first started appearing in Chinese records during the Han dynasty as the k'un-lun po or kunlun bo ("ship of the k'un-lun [dark-skinned southern people]"). These ships used two types of sail of their invention, the junk sail and tanja sail. Large ships are about 50–60 metres (164–197 ft) long, had tall freeboard, each carrying provisions enough for a year, and could carry 200–1000 people. The Chinese recorded that these Southeast Asian ships were hired for passage to South Asia by Chinese Buddhist pilgrims and travelers, because they did not build seaworthy ships of their own until around the 8–9th century AD. Austronesians (especially from western Island Southeast Asia) were trading in the Indian Ocean as far as Africa during this period. By around 50 to 500 AD, a group of Austronesians, believed to be from the southeastern coasts of Borneo (possibly a mixed group related to the modern Ma'anyan, Banjar, and/or the Dayak people) crossed the Indian Ocean and colonized Madagascar. This resulted in the introduction of outrigger canoe technology to non-Austronesian cultures in the East African coast. China The ancient Chinese also built fluvial ramming vessels as in the Greco-Roman tradition of the trireme, although oar-steered ships in China lost favor very early on since it was in the 1st century China that the stern-mounted rudder was first developed. This was dually met with the introduction of the Han dynasty junk ship design in the same century. The Chinese were using square sails during the Han dynasty and adopted the Austronesian junk sail later in the 12th century. Iconographic remains show that Chinese ships before the 12th century used square sails, and the junk rig of Chinese ships is believed to be developed from tilted sails. Southern Chinese junks were based on keeled and multi-planked Austronesian ship known as po by the Chinese, from the Old Javanese parahu, Javanese prau, or Malay perahu – large ship. Southern Chinese junks showed characteristics of Austronesian ships that they are made using timbers of tropical origin, with keeled, V-shaped hull. This is different from northern Chinese junks, which are developed from flat-bottomed riverine boats. The northern Chinese junks were primarily built of pine or fir wood, had flat bottoms with no keel, water-tight bulkheads with no frames, transom (squared) stern and stem, and have their planks fastened with iron nails or clamps. It was unknown when the Chinese people started adopting Southeast Asian (Austronesian) shipbuilding techniques. They may have been started as early as the 8th century, but the development was gradual and the true ocean-going Chinese junks did not appear suddenly. The word "po" survived in Chinese long after, referring to the large ocean-going junks. Mediterranean In September 2011, archeological investigations done at the site of Portus in Rome revealed inscriptions in a shipyard constructed during the reign of Trajan (98–117) that indicated the existence of a shipbuilders guild. Early 2nd millennium AD Austronesia Roughly at this time is the last migration wave of the Austronesian expansion, when the Polynesian islands spread over vast distances across the Pacific Ocean were being colonized by the (Austronesian) Polynesians from Island Melanesia using double-hulled voyaging catamarans. At its furthest extent, there is a possibility that they may have reached the Americas. After the 11th century, a new type of ship called djong or jong was recorded in Java and Bali. This type of ship was built using wooden dowels and treenails, unlike the kunlun bo which used vegetal fibres for lashings. The empire of Majapahit used jong, built in northern Java, for transporting troops overseas. The jongs were transport ships which could carry 100–2000 tons of cargo and 50–1000 people, 28.99–88.56 meter in length. The exact number of jong fielded by Majapahit is unknown, but the largest number of jong deployed in an expedition is about 400 jongs, when Majapahit attacked Pasai, in 1350. Europe Until recently, Viking longships were seen as marking an advance on traditional clinker-built hulls where leather thongs were used to join plank boards. This consensus has recently been challenged. Haywood has argued that earlier Frankish and Anglo-Saxon nautical practice was much more accomplished than had been thought and has described the distribution of clinker vs. carvel construction in Western Europe (see map ). An insight into shipbuilding in the North Sea/Baltic areas of the early medieval period was found at Sutton Hoo, England, where a ship was buried with a chieftain. The ship was long and wide. Upward from the keel, the hull was made by overlapping nine strakes on either side with rivets fastening the oaken planks together. It could hold upwards of thirty men. Sometime around the 12th century, northern European ships began to be built with a straight sternpost, enabling the mounting of a rudder, which was much more durable than a steering oar held over the side. Development in the Middle Ages favored "round ships", with a broad beam and heavily curved at both ends. Another important ship type was the galley, which was constructed with both sails and oars. The first extant treatise on shipbuilding was written by Michael of Rhodes, a man who began his career as an oarsman on a Venetian galley in 1401 and worked his way up into officer positions. He wrote and illustrated a book that contains a treatise on shipbuilding, a treatise on mathematics, much material on astrology, and other materials. His treatise on shipbuilding treats three kinds of galleys and two kinds of round ships. China Shipbuilders in the Ming dynasty (1368~1644) were not the same as the shipbuilders in other Chinese dynasties, due to hundreds of years of accumulated experiences and rapid changes in the Ming dynasty. Shipbuilders in the Ming dynasty primarily worked for the government, under command of the Ministry of Public Works. During the early years of the Ming dynasty, the Ming government maintained an open policy towards sailing. Between 1405 and 1433, the government conducted seven diplomatic Ming treasure voyages to over thirty countries in Southeast Asia, India, the Middle East and Eastern Africa. The voyages were initiated by the Yongle Emperor, and led by the Admiral Zheng He. Six voyages were conducted under the Yongle Emperor's reign, the last of which returned to China in 1422. After the Yongle Emperor's death in 1424, his successor the Hongxi Emperor ordered the suspension of the voyages. The seventh and final voyage began in 1430, sent by the Xuande Emperor. Although the Hongxi and Xuande Emperors did not emphasize sailing as much as the Yongle Emperor, they were not against it. This led to a high degree of commercialization and an increase in trade. Large numbers of ships were built to meet the demand. The Ming voyages were large in size, numbering as many as 300 ships and 28,000 men. The shipbuilders were brought from different places in China to the shipyard in Nanjing, including Zhejiang, Jiangxi, Fujian, and Huguang (now the provinces of Hubei and Hunan). One of the most famous shipyards was Long Jiang Shipyard (:zh:龙江船厂), located in Nanjing near the Treasure Shipyard where the ocean-going ships were built. The shipbuilders could build 24 models of ships of varying sizes. Several types of ships were built for the voyages, including Shachuan (沙船), Fuchuan (福船) and Baochuan (treasure ship) (宝船). Zheng He's treasure ships were regarded as Shachuan types, mainly because they were made in the treasure shipyard in Nanjing. Shachuan, or 'sand-ships', are ships used primarily for inland transport. However, in recent years, some researchers agree that the treasure ships were more of the Fuchuan type. It is said in vol. 176 of San Guo Bei Meng Hui Bian (三朝北盟汇编) that ships made in Fujian are the best ones. Therefore, the best shipbuilders and laborers were brought from these places to support Zheng He's expedition. The shipyard was under the command of Ministry of Public Works. The shipbuilders had no control over their lives. The builders, commoner's doctors, cooks and errands had lowest social status. The shipbuilders were forced to move away from their hometown to the shipyards. There were two major ways to enter the shipbuilder occupation: family tradition, or apprenticeship. If a shipbuilder entered the occupation due to family tradition, the shipbuilder learned the techniques of shipbuilding from his family and is very likely to earn a higher status in the shipyard. Additionally, the shipbuilder had access to business networking that could help to find clients. If a shipbuilder entered the occupation through an apprenticeship, the shipbuilder was likely a farmer before he was hired as a shipbuilder, or he was previously an experienced shipbuilder. Many shipbuilders working in the shipyard were forced into the occupation. The ships built for Zheng He's voyages needed to be waterproof, solid, safe, and have ample room to carry large amounts of trading goods. Therefore, due to the highly commercialized society that was being encouraged by the expeditions, trades, and government policies, the shipbuilders needed to acquire the skills to build ships that fulfil these requirements. Shipbuilding was not the sole industry utilising Chinese lumber at that time; the new capital was being built in Beijing from approximately 1407 onwards, which required huge amounts of high-quality wood. These two ambitious projects commissioned by Emperor Yongle would have had enormous environmental and economic effects, even if the ships were half the dimensions given in the History of Ming. Considerable pressure would also have been placed on the infrastructure required to transport the trees from their point of origin to the shipyards. Shipbuilders were usually divided into different groups and had separate jobs. Some were responsible for fixing old ships; some were responsible for making the keel and some were responsible for building the helm. It was the keel that determined the shape and the structure of the hull of Fuchuan Ships. The keel is the middle of the bottom of the hull, constructed by connecting three sections; stern keel, main keel and poop keel. The hull spreads in the arc towards both sides forming the keel. The helm was the device that controls direction when sailing. It was a critical invention in shipbuilding technique in ancient China and was only used by the Chinese for a fairly long time. With a developing recognition of its function, the shape and configuration of the helm was continually improved by shipbuilders. The shipbuilders not only needed to build the ship according to design, but needed to acquire the skills to improve the ships. After 1477, the Ming government reversed its open maritime policies, enacting a series of isolationist policies in response to piracy. The policies, called Haijin (sea ban), lasted until the end of the Ming dynasty in 1644. During this period, Chinese navigation technology did not make any progress and even declined in some aspect. Indian Ocean In the Islamic world, shipbuilding thrived at Basra and Alexandria. The dhow, felucca, baghlah, and the sambuk became symbols of successful maritime trade around the Indian Ocean from the ports of East Africa to Southeast Asia and the ports of Sindh and Hind (India) during the Abbasid period. Early modern Bengal Mughal Empire had a large shipbuilding industry, which was largely centred in the Bengal Subah. Economic historian Indrajit Ray estimates shipbuilding output of Bengal during the sixteenth and seventeenth centuries at 223,250 tons annually, compared with 23,061 tons produced in nineteen colonies in North America from 1769 to 1771. He also assesses ship repairing as very advanced in Bengal. West Africa Documents from 1506, for example, refer to watercraft on the Sierra Leone river carrying 120 men. Others refer to Guinea coast peoples using war canoes of varying sizes – some 70 feet in length, 7–8 feet broad, with sharp pointed ends, rowing benches on the side, and quarterdecks or forecastles build of reeds. The watercraft included miscellaneous facilities, such as cooking hearths, and storage spaces for the crew's sleeping mats. From the 17th century, some kingdoms added brass or iron cannons to their vessels. By the 18th century, however, the use of swivel cannons on war canoes accelerated. The city-state of Lagos, for instance, deployed war canoes armed with swivel cannons. Europe With the development of the carrack, the west moved into a new era of ship construction by building the first regular oceangoing vessels. In a relatively short time, these ships grew to an unprecedented size, complexity, and cost. Shipyards became large industrial complexes, and the ships built were financed by consortia of investors. These considerations led to the documentation of design and construction practices in what had previously been a secretive trade run by master shipwrights and ultimately led to the field of naval architecture, in which professional designers and draftsmen played an increasingly important role. Even so, construction techniques changed only very gradually. The ships of the Napoleonic Wars were still built more or less to the same basic plan as those of the Spanish Armada of two centuries earlier, although there had been numerous subtle improvements in ship design and construction throughout this period. For instance, the introduction of tumblehome, adjustments to the shapes of sails and hulls, the introduction of the wheel, the introduction of hardened copper fastenings below the waterline, the introduction of copper sheathing as a deterrent to shipworm and fouling, etc. Industrial Revolution Initially copying wooden construction traditions with a frame over which the hull was fastened, Isambard Kingdom Brunel's of 1843 was the first radical new design, being built entirely of wrought iron. Despite her success, and the great savings in cost and space provided by the iron hull, compared to a copper-sheathed counterpart, there remained problems with fouling due to the adherence of weeds and barnacles. As a result, composite construction remained the dominant approach where fast ships were required, with wooden timbers laid over an iron frame (Cutty Sark is a famous example). Later Great Britains iron hull was sheathed in wood to enable it to carry a copper-based sheathing. Brunel's represented the next great development in shipbuilding. Built-in association with John Scott Russell, it used longitudinal stringers for strength, inner and outer hulls, and bulkheads to form multiple watertight compartments. Steel also supplanted wrought iron when it became readily available in the latter half of the 19th century, providing great savings when compared with iron in cost and weight. Wood continued to be favored for the decks. During World War II, the need for cargo ships was so great that construction time for Liberty ships went from initially eight months or longer, down to weeks or even days. They employed production line and prefabrication techniques such as those used in shipyards today. The total number of dry-cargo ships built in the United States in a 15-year period just before the war was a grand total of two. During the war, thousands of Liberty ships and Victory ships were built, many of them in shipyards that did not exist before the war. And, they were built by a workforce consisting largely of women and other inexperienced workers who had never seen a ship before (or even the ocean). Worldwide shipbuilding industry After World War II, shipbuilding (which encompasses the shipyards, the marine equipment manufacturers, and many related service and knowledge providers) grew as an important and strategic industry in a number of countries around the world. This importance stems from: The large number of skilled workers required directly by the shipyard, along with supporting industries such as steel mills, railroads and engine manufacturers; and A nation's need to manufacture and repair its own navy and vessels that support its primary industries Historically, the industry has suffered from the absence of global rules and a tendency towards (state-supported) over-investment due to the fact that shipyards offer a wide range of technologies, employ a significant number of workers, and generate income as the shipbuilding market is global. Japan used shipbuilding in the 1950s and 1960s to rebuild its industrial structure; South Korea started to make shipbuilding a strategic industry in the 1970s, and China is now in the process of repeating these models with large state-supported investments in this industry. Conversely, Croatia is privatising its shipbuilding industry. As a result, the world shipbuilding market suffers from over-capacities, depressed prices (although the industry experienced a price increase in the period 2003–2005 due to strong demand for new ships which was in excess of actual cost increases), low profit margins, trade distortions and widespread subsidisation. All efforts to address the problems in the OECD have so far failed, with the 1994 international shipbuilding agreement never entering into force and the 2003–2005 round of negotiations being paused in September 2005 after no agreement was possible. After numerous efforts to restart the negotiations these were formally terminated in December 2010. The OECD's Council Working Party on Shipbuilding (WP6) will continue its efforts to identify and progressively reduce factors that distort the shipbuilding market. Where state subsidies have been removed and domestic industrial policies do not provide support in high labor cost countries, shipbuilding has gone into decline. The British shipbuilding industry is a prime example of this with its industries suffering badly from the 1960s. In the early 1970s British yards still had the capacity to build all types and sizes of merchant ships but today they have been reduced to a small number specialising in defence contracts, luxury yachts and repair work. Decline has also occurred in other European countries, although to some extent this has reduced by protective measures and industrial support policies. In the US, the Jones Act (which places restrictions on the ships that can be used for moving domestic cargoes) has meant that merchant shipbuilding has continued, albeit at a reduced rate, but such protection has failed to penalise shipbuilding inefficiencies. The consequence of this is that contract prices are far higher than those of any other country building oceangoing ships. Present day shipbuilding Beyond the 2000s, the three East Asian manufacturing powerhouses, China, South Korea and Japan, have dominated world shipbuilding by completed gross tonnage. China State Shipbuilding Corporation, China Shipbuilding Industry Corporation, Hyundai Heavy Industries, Samsung Heavy Industries, Daewoo Shipbuilding & Marine Engineering and Imabari Shipbuilding supply most of the global market for large container, bulk carrier, tanker and Ro-ro ships. During the early 2020s, Chinese shipbuilders saw an increase in orders as operators demand greener fleets. , China’s shipbuilding output, newly received orders and orders-on-hand accounted for 50%, 67%, and 55% of the global market share, respectively, with double-digital growth of all three indexes compared to the previous year. When referring to the ship type, then China, South Korea and Japan are the producing countries of the carrier ships as mentioned above. While Italy, France, Finland, Germany, United Kingdom and other European countries are the makers of cruise ships (the most), icebreakers, crane vessel and so on. The market share of European ship builders began to decline in the 1960s as they lost work to Japan in the same way Japan most recently lost their work to South Korea and China. Over the four years from 2007, the total number of employees in the European shipbuilding industry declined from 150,000 to 115,000. In 2022, some key shipbuilders in Europe are Fincantieri, Damen Group, Navantia, Naval Group and BAE Systems. Shipbuilding output of the United States also underwent a similar change. The US is ranked the 10th largest shipbuilder worldwide. The top companies that build large naval vessels, such as aircraft carriers and cruisers, include Huntington Ingalls, Bollinger and General Dynamics. In the small to medium military vessels category, key shipbuilders include Vigor Industrial, and VT Halter Marine. As the US Navy is shifting to a new fleet architecture that is more widely distributed, unmanned surface vehicles (USVs) development is rapidly propelled to higher priority. Key strategic Program of Record includes prototyping and construction of up to 9 MUSVs, for which a sole contract was awarded to L3Harris Technologies, who partnered with Swiftships to build the MUSVs. In 2018, the US Defense Department initiated Overlord Program, and developed USV Prototypes 1 (NOMAD) and 2 (RANGER). Both of them took part in multiple fleet level exercises and demonstrations, traveled in autonomous mode, and tested numerous payloads. Nomam, formerly known as Riley Claire, is a converted offshore patrol vessel, which was built by Swiftships. The objective of the Ghost Fleet Overlord program is to convert large, commercial vessels to autonomous systems. , China builds nearly 75% of new vessels on order and has the capacity to build more ships in one month than the United States builds in a year. Modern shipbuilding manufacturing techniques Modern shipbuilding makes considerable use of prefabricated sections. Entire multi-deck segments of the hull or superstructure will be built elsewhere in the yard, transported to the building dock or slipway, then lifted into place. This is known as "block construction". The most modern shipyards pre-install equipment, pipes, electrical cables, and any other components within the blocks, to minimize the effort needed to assemble or install components deep within the hull once it is welded together. Ship design work, also called naval architecture, may be conducted using a ship model basin. Previously, loftsmen at the mould lofts of shipyards were responsible for taking the dimensions, and details from drawings and plans and translating this information into templates, battens, ordinates, cutting sketches, profiles, margins and other data. However, since the early 1970s computer-aided design became normal for the shipbuilding design and lofting process. Modern ships, since roughly 1940, have been produced almost exclusively of welded steel. Early welded steel ships used steels with inadequate fracture toughness, which resulted in some ships suffering catastrophic brittle fracture structural cracks (see problems of the Liberty ship). Since roughly 1950, specialized steels such as ABS Steels with good properties for ship construction have been used. Although it is commonly accepted that modern steel has eliminated brittle fracture in ships, some controversy still exists. Brittle fracture of modern vessels continues to occur from time to time because grade A and grade B steel of unknown toughness or fracture appearance transition temperature (FATT) in ships' side shells can be less than adequate for all ambient conditions. As modern shipbuilding panels on a panel line become lighter and thinner, the laser hybrid welding technique is utilized. The laser hybrid blend focuses a higher energy beam on the material to be joined, allowing it to keyhole with a much higher depth to width ratio than comparative traditional welding techniques. Typically a MIG process trails the keyhole providing filler material for the weld joint. This allows for very high penetration without excessive heat input from decreased weld metal deposited leading to less distortion and welding at higher travel speeds. Ship repair industry All ships need repair work at some point in their working lives. A part of these jobs must be carried out under the supervision of the classification society. A lot of maintenance is carried out while at sea or in port by ship's crew. However, a large number of repair and maintenance works can only be carried out while the ship is out of commercial operation, in a ship repair yard. Prior to undergoing repairs, a tanker must dock at a deballasting station for completing the tank cleaning operations and pumping ashore its slops (dirty cleaning water and hydrocarbon residues).	Technology	Maritime transport	null
187386	https://en.wikipedia.org/wiki/Egyptian%20vulture	Egyptian vulture	The Egyptian vulture (Neophron percnopterus), also called the white scavenger vulture or pharaoh's chicken, is a small Old World vulture in the monotypic genus Neophron. It is widely distributed from the Iberian Peninsula, North Africa, West Asia and India. The contrasting underwing pattern and wedge-shaped tail make it distinctive in flight as it soars in thermals during the warmer parts of the day. Egyptian vultures feed mainly on carrion but are opportunistic and will prey on small mammals, birds, and reptiles. They also feed on the eggs of other birds, breaking larger ones by tossing a large pebble onto them. The use of tools is rare in birds and apart from the use of a pebble as a hammer, Egyptian vultures also use twigs to roll up wool for use in their nest. Egyptian vultures that breed in the temperate regions migrate south in winter while tropical populations are relatively sedentary. Populations of this species declined in the 20th century and some island populations are endangered by hunting, accidental poisoning, and collision with power lines. Taxonomy and systematics The Egyptian vulture was formally described by the Swedish naturalist Carl Linnaeus in 1758 in the tenth edition of his Systema Naturae under the binomial name Vultur percnopterus. The genus Neophron was created by Jules-César Savigny in the first natural history volume of the Description de l'Égypte (1809). The genus Neophron contains only a single extant species. A few prehistoric species from the Neogene period in North America placed in the genus Neophrontops (the name meaning "looks like Neophron") are believed to have been very similar to these vultures in lifestyle, but the genetic relationships are unclear. A fossil species Neophron lolis has been described from the late Miocene of Spain. The genus Neophron is considered to represent the oldest branch of the vultures which consists of separated (or polyphyletic) clades. Along with its nearest evolutionary relatives, the lammergeier (Gypaetus barbatus) and the palm-nut vulture (Gypohierax angolensis), they are sometimes placed in a separate subfamily, the Gypaetinae. Subspecies There are three widely recognised subspecies of the Egyptian vulture, although there is considerable gradation due to movement and intermixing of the populations. The nominate subspecies, N. p. percnopterus, with a dark grey bill, has the largest range, occurring in southern Europe, northern Africa, the Middle East, Central Asia, and north-western India. Populations breeding in the temperate zone migrate south during winter. The Indian subcontinent is the range of subspecies N. p. ginginianus, the smallest of the three subspecies, which is identifiable by a pale yellow bill. The subspecies name is derived from Gingee in southern India from where the French explorer Pierre Sonnerat described it as Le Vautour de Gingi and it was given a Latin name by John Latham in his Index Ornithologicus (1790). A small population that is found only in the eastern Canary Islands was found to be genetically distinct and identified as a new subspecies, N. p. majorensis in 2002. Known locally as the guirre they are genetically more distant from N. p. percnopterus, significantly greater even than N. p. ginginianus is from N. p. percnopterus. Unlike neighbouring populations in Africa and southern Europe, it is non-migratory and consistently larger in size. The subspecies name majorensis is derived from "Majorata", the ancient name for the island of Fuerteventura. The island was named by Spanish conquerors in the 15th century after the "Majos", the main native Guanche tribe there. One study in 2010 suggested that the species established on the island about 2,500 years ago when the island was first colonized by humans. Nikolai Zarudny and Härms described a subspecies, rubripersonatus, from Baluchistan in 1902. This was described as having a deeper reddish orange skin on the head and a yellow-tipped dark bill. This has rarely been considered a valid subspecies but the intermediate pattern of bill colouration suggests intermixing of subspecies. Etymology The genus name is derived from Greek mythology. Timandra was the mother of Neophron. Aegypius was a friend of Neophron and about the same age. It upset Neophron to know that his mother Timandra was having a love affair with Aegypius. Seeking revenge, Neophron made advances towards Aegypius' mother, Bulis. Neophron succeeded and enticed Bulis into entering the dark chamber where his mother and Aegypius were to meet soon. Neophron then distracted his mother, tricking Aegypius into entering the chamber and sleeping with his own mother Bulis. When Bulis discovered the deception she gouged out the eyes of her son Aegypius before killing herself. Aegypius prayed for revenge and Zeus, on hearing the prayer, changed Aegypius and Neophron into vultures. "Percnopterus" is derived from Greek for "black wings": "περκνός" (perknos, meaning "blue-black") and πτερόν (pteron, meaning wing). Description The adult's plumage is white, with black flight feathers in the wings. Wild birds usually appear soiled with a rusty or brown shade to the white plumage, derived from mud or iron-rich soil. Captive specimens without access to soil have clean white plumage. It has been suggested as a case of cosmetic colouration. The bill is slender and long, and the tip of the upper mandible is hooked. The nostril is an elongated horizontal slit. The neck feathers are long and form a hackle. The wings are pointed, with the third primary being the longest; the tail is wedge shaped. The legs are pink in adults and grey in juveniles. The claws are long and straight, and the third and fourth toes are slightly webbed at the base. The bill is black in the nominate subspecies but pale or yellowish in adults of the smaller Indian ginginianus. Rasmussen and Anderton (2005) suggest that this variation may need further study, particularly due to the intermediate black-tipped bill described in rubripersonatus. The facial skin is yellow and unfeathered down to the throat. The sexes are indistinguishable in plumage but breeding males have a deeper orange facial skin colour than females. Females average slightly larger and are about 10–15% heavier than males. Young birds are blackish or chocolate brown with black and white patches. The adult plumage is attained only after about five years. The adult Egyptian vulture measures from the point of the beak to the extremity of the tail feathers. In the smaller N. p. ginginianus males are about long while females are long. The wingspan is about 2.7 times the body length. Birds from Spain weigh about while birds of the Canary Island subspecies majorensis, representing a case of island gigantism, are heavier with an average weight of . The Egyptian vulture is one of the smallest true Old World vulture, the only smaller species appears to be the marginally lighter palm-nut vulture (which may be an outlier from other vultures). Additionally, the hooded vulture is only scarcely larger than the Egyptian species. Distribution and movements Egyptian vultures are widely distributed across the Old World with their breeding range from southern Europe to northern Africa east to western and southern Asia. They are rare vagrants in Sri Lanka. They occur mainly on the dry plains and lower hills. In the Himalayas, they go up to about in summer. In Armenia, breeding pairs have been found up to 2,300 meters a.s.l. Most Egyptian vultures in the subtropical zone of Europe migrate south to Africa in winter. Vagrants may occur as far south as in South Africa although they bred in the Transkei region prior to 1923. They nest mainly on rocky cliffs, sometimes adopting ledges on tall buildings in cities and on large trees. Like many other large soaring migrants, they avoid making long crossings over water. Italian birds cross over through Sicily and into Tunisia making short sea crossings by passing through the islands of Marettimo and Pantelleria with rare stops on the island country of Malta. Those that migrate through the Iberian Peninsula cross into Africa over the Strait of Gibraltar while others cross further east through the Levant. In summer, some African birds fly further north into Europe and vagrants have been recorded in England, Ireland, and southern Sweden. Migrating birds can sometimes cover in a single day until they reach the southern edge of the Sahara, from their summer home. Young birds that have not reached breeding age may overwinter in the grassland and semi-desert regions of the Sahel. Fossil record Fossils of the Egyptian vulture found in the Nefud Desert of Saudi Arabia are estimated to date to the Middle Pleistocene about 500,000 years ago. Behaviour and ecology The Egyptian vulture is usually seen singly or in pairs, soaring in thermals along with other scavengers and birds of prey, or perched on the ground or atop a building. On the ground, they walk with a waddling gait. They feed on a range of food, including mammal faeces (including those of humans), insects in dung, carrion, vegetable matter, and sometimes small animals. It is the only Old World vulture species that regularly feeds on faeces. The carotenoids (primarily lutein) that the vultures absorb from the vegetal matter in the excrement that they ingest results in their bright yellow face colouration. When it joins other vulture species at a dead animal, it tends to stay on the periphery and waits until the larger species leave. Pairs may also scrounge for food from other vultures, particularly griffons. Recently fledged young will sometimes fly to other nests, competing with young vultures for food, stealing or even soliciting food from the (unrelated) adults bringing food. Wild rabbits (Oryctolagus cuniculus) form a significant part of the diet of Spanish vultures. In the Iberian Peninsula, landfills are an important food source, with the vultures more likely to occupy territories close to landfill sites. Studies suggest that they feed on ungulate faeces to obtain carotenoid pigments responsible for their bright yellow and orange facial skin. The ability to assimilate carotenoid pigments may serve as a reliable signal of fitness. Egyptian vultures are mostly silent but make high-pitched mewing or hissing notes at the nest and screeching noises when squabbling at a carcass. Young birds have been heard making a hissing croak in flight. They also hiss or growl when threatened or angry.Egyptian vultures roost communally on large trees, buildings or on cliffs. Roost sites are usually chosen close to a dump site or other suitable foraging area. In Spain and Morocco, summer roosts are formed mainly by immature birds. The favourite roost trees tended to be large dead pines. The number of adults at the roost increases towards June. It is thought that breeding adults may be able to forage more efficiently by joining the roost and following others to the best feeding areas. Breeding birds that failed to raise young may also join the non-breeding birds at the roost during June. Allopreening has been observed in Canarian Egyptian vultures between mated pairs of individuals as well as pairs of unrelated and same-sex individuals, particularly females. Breeding The breeding season is in spring. During the beginning of the breeding season, courting pairs soar high together and one or both may make steep spiralling or swooping dives. The birds are monogamous and pair bonds may be maintained for more than one breeding season and the same nest sites may be reused each year. The nest is an untidy platform of twigs lined with rags and placed on a cliff ledge, building, or the fork of a large tree. Old nest platforms of eagles may also be taken over. Nests placed on the ground are rare but have been recorded in subspecies N. p. ginginianus and N. p. majorensis. Extra-pair copulation with neighbouring birds has been recorded and may be a reason for adult males to stay close to the female before and during the egg laying period. Females may sometimes associate with two males and all three help in raising the brood. The typical clutch consists of two eggs which are incubated in turns by both parents. The eggs are brick red with the broad end covered more densely with blotches of red, brown, and black. The parents begin incubating soon after the first egg is laid leading to asynchronous hatching. The first egg hatches after about 42 days. The second chick may hatch three to five days later and a longer delay increases the likelihood that it will die of starvation. In cliffs where the nests are located close to each other, young birds have been known to clamber over to neighbouring nests to obtain food. In the Spanish population, young fledge and leave the nest after 90 to 110 days. Fledged birds continue to remain dependent on their parents for at least a month. Once the birds begin to forage on their own, they move away from their parents' territory; young birds have been found nearly 500 km away from their nest site. One-year-old European birds migrate to Africa and stay there for at least one year. A vulture that fledged in France stayed in Africa for three years before migrating north in spring. After migrating back to their breeding areas, young birds move widely in search of good feeding territories and mates. The full adult plumage is attained in the fourth or fifth year. Egyptian vultures have been known to live for up to 37 years in captivity and at least 21 years in the wild. The probability of survival in the wild varies with age, increasing till the age of 2 and then falling at the age of 5. Older birds have an annual survival probability varying from 0.75 for non-breeders to 0.83 for breeding birds. Tool use The nominate population, especially in Africa, is known for its use of stones as tools. When a large egg, such as that of an ostrich or bustard, is located, the bird walks up to it with a large pebble held in its bill and tosses the pebble by swinging the neck down over the egg. The operation is repeated until the egg cracks from the blows. They prefer using rounded pebbles to jagged rocks. This behaviour, although believed to have been first reported by Jane Goodall in 1966, was actually already known to Africans and was first reported by J. G. Wood in 1877. However, this has only been reported in Africa and has not been recorded in N. p. ginginianus. Tests with both hand-reared and wild birds suggest that the behaviour is innate, not learnt by observing other birds, and elicited once they associate eggs with food and have access to pebbles. Their ability to deal with ostrich eggs is utilized by brown-necked ravens which form groups that wait for the eggs to be broken before collectively mobbing the vultures and engaging in kleptoparasitism. Another case of tool-use described from Bulgaria involves the use of a twig as a tool to roll up and gather strands of wool to use for lining the nest. Threats and conservation Healthy adults do not have many predators, but human activities pose many threats. Collisions with power lines, hunting, intentional poisoning, lead accumulation from ingesting gunshot in carcasses, and pesticide accumulation take a toll on populations. Young birds at the nest are sometimes taken by golden eagles, eagle owls, and red foxes. Only rarely do adult birds attempt to drive away predators. Young birds that fall off of cliff ledges may be preyed on by mammalian predators such as jackals, foxes and wolves. Like all birds they serve as hosts for ectoparasitic birdlice including Aegypoecus perspicuus as well as organisms that live within them such as mycoplasmas. Egyptian vulture populations have declined in most parts of its range. In Europe and most of the Middle East, populations in 2001 were half of those from 1980. In India, the decline has been rapid with a 35% decrease each year since 1999. In 1967–70, the area around Delhi was estimated to have 12,000–15,000 of these vultures, with an average density of about 5 pairs per 10 km2. The exact cause of the decline is not known, but has been linked with the use of the NSAID Diclofenac, which has been known to cause death in Gyps vultures. In Italy, the number of breeding pairs declined from 30 in 1970 to 9 in the 1990s. Nearly all breeding failures were due to human activities. In Spain, which holds about 50% of the European population suggested causes of decline include poisoning by accumulation of lead, pesticides (especially due to large-scale use in the control of Schistocerca gregaria locust swarms), and electrocution. Windfarms may also pose a threat. Poorly designed power transmission lines in east Africa electrocute many wintering vultures. A shortage of carrion resulting from new rules for disposal of dead animals following the outbreak of Bovine Spongiform Encephalitis in parts of Europe during 2000 may have also had an effect on some populations. In Armenia direct persecution for trophy and for local illegal trade of animals as pets has been recorded. The population of Egyptian vultures in the Canary Islands has been isolated from those in Europe and Africa for a significant period of time leading to genetic differentiation. The vulture population there declined by 30% in the ten years between 1987 and 1998. The Canarian Egyptian vulture was historically common, occurring on the islands of La Gomera, Tenerife, Gran Canaria, Fuerteventura, and Lanzarote. It is now restricted to Fuerteventura and Lanzarote, the two easternmost islands. The total population in 2000 was estimated at 130 individuals, including 25–30 breeding pairs. Island birds also appear to accumulate significant amounts of lead from scavenging on hunted animal carcasses. The long-term effect of this poison at a sub-lethal level is not known, though it is known to alter the mineralization of their bones. In order to provide safe and uncontaminated food for nesting birds, attempts have been made to create "vulture restaurants" where carcasses are made available. However, these interventions may also encourage other opportunist predators and scavengers to concentrate at the site and pose a threat to vultures nesting in the vicinity. Since 2012, conservation efforts have been implemented to protect Egyptian vultures along breeding grounds, migration routes and wintering areas of the eastern European population; these measures include monitoring, nest guarding, supplementary feeding, insulating hazardous electric power lines and removing poison baits and carcasses with trained dogs. Adult annual survival and juvenile monthly survival appear to have increased, leading to a notable rise in population growth. In culture The Bible makes a reference to the Egyptian vulture under the Hebrew name of rachamah/racham which has been translated into English as "gier-eagle". In Ancient Egypt, several hieroglyphs include the Egyptian vulture including what is listed as G1 in the Gardiner's sign list - . The bird was held sacred to Isis and Mut in ancient Egyptian religion. The use of the vulture as a symbol of royalty in Egyptian culture and their protection by Pharaonic law made the species common on the streets of Egypt and gave rise to the name "pharaoh's chicken". The habit of coprophagy in Egyptian vultures gives them the Spanish names of "churretero" and "moñiguero", which mean "dung-eater". British sportsmen in colonial India considered them to be among the ugliest birds, and their habit of feeding on faeces was particularly despised. In British India they were known as "shawks" a contraction of shit-hawk. A southern Indian temple at Thirukalukundram near Chengalpattu was famed for a pair of birds that reputedly visited the temple for "centuries". These birds were ceremonially fed by the temple priests and arrived before noon to feed on offerings made from rice, wheat, ghee, and sugar. Although normally punctual, the failure of the birds to turn up was attributed to the presence of "sinners" among the onlookers. Legend has it the vultures (or "eagles") represented eight sages who were punished by Shiva, with two of them leaving in each of a series of epochs.	Biology and health sciences	Accipitrimorphae	Animals
187446	https://en.wikipedia.org/wiki/Orientability	Orientability	In mathematics, orientability is a property of some topological spaces such as real vector spaces, Euclidean spaces, surfaces, and more generally manifolds that allows a consistent definition of "clockwise" and "anticlockwise". A space is orientable if such a consistent definition exists. In this case, there are two possible definitions, and a choice between them is an orientation of the space. Real vector spaces, Euclidean spaces, and spheres are orientable. A space is non-orientable if "clockwise" is changed into "counterclockwise" after running through some loops in it, and coming back to the starting point. This means that a geometric shape, such as , that moves continuously along such a loop is changed into its own mirror image . A Möbius strip is an example of a non-orientable space. Various equivalent formulations of orientability can be given, depending on the desired application and level of generality. Formulations applicable to general topological manifolds often employ methods of homology theory, whereas for differentiable manifolds more structure is present, allowing a formulation in terms of differential forms. A generalization of the notion of orientability of a space is that of orientability of a family of spaces parameterized by some other space (a fiber bundle) for which an orientation must be selected in each of the spaces which varies continuously with respect to changes in the parameter values. Orientable surfaces A surface S in the Euclidean space R3 is orientable if a chiral two-dimensional figure (for example, ) cannot be moved around the surface and back to where it started so that it looks like its own mirror image (). Otherwise the surface is non-orientable. An abstract surface (i.e., a two-dimensional manifold) is orientable if a consistent concept of clockwise rotation can be defined on the surface in a continuous manner. That is to say that a loop going around one way on the surface can never be continuously deformed (without overlapping itself) to a loop going around the opposite way. This turns out to be equivalent to the question of whether the surface contains no subset that is homeomorphic to the Möbius strip. Thus, for surfaces, the Möbius strip may be considered the source of all non-orientability. For an orientable surface, a consistent choice of "clockwise" (as opposed to counter-clockwise) is called an orientation, and the surface is called oriented. For surfaces embedded in Euclidean space, an orientation is specified by the choice of a continuously varying surface normal n at every point. If such a normal exists at all, then there are always two ways to select it: n or −n. More generally, an orientable surface admits exactly two orientations, and the distinction between an oriented surface and an orientable surface is subtle and frequently blurred. An orientable surface is an abstract surface that admits an orientation, while an oriented surface is a surface that is abstractly orientable, and has the additional datum of a choice of one of the two possible orientations. Examples Most surfaces encountered in the physical world are orientable. Spheres, planes, and tori are orientable, for example. But Möbius strips, real projective planes, and Klein bottles are non-orientable. They, as visualized in 3-dimensions, all have just one side. The real projective plane and Klein bottle cannot be embedded in R3, only immersed with nice intersections. Note that locally an embedded surface always has two sides, so a near-sighted ant crawling on a one-sided surface would think there is an "other side". The essence of one-sidedness is that the ant can crawl from one side of the surface to the "other" without going through the surface or flipping over an edge, but simply by crawling far enough. In general, the property of being orientable is not equivalent to being two-sided; however, this holds when the ambient space (such as R3 above) is orientable. For example, a torus embedded in can be one-sided, and a Klein bottle in the same space can be two-sided; here refers to the Klein bottle. Orientation by triangulation Any surface has a triangulation: a decomposition into triangles such that each edge on a triangle is glued to at most one other edge. Each triangle is oriented by choosing a direction around the perimeter of the triangle, associating a direction to each edge of the triangle. If this is done in such a way that, when glued together, neighboring edges are pointing in the opposite direction, then this determines an orientation of the surface. Such a choice is only possible if the surface is orientable, and in this case there are exactly two different orientations. If the figure can be consistently positioned at all points of the surface without turning into its mirror image, then this will induce an orientation in the above sense on each of the triangles of the triangulation by selecting the direction of each of the triangles based on the order red-green-blue of colors of any of the figures in the interior of the triangle. This approach generalizes to any n-manifold having a triangulation. However, some 4-manifolds do not have a triangulation, and in general for n > 4 some n-manifolds have triangulations that are inequivalent. Orientability and homology If H1(S) denotes the first homology group of a closed surface S, then S is orientable if and only if H1(S) has a trivial torsion subgroup. More precisely, if S is orientable then H1(S) is a free abelian group, and if not then H1(S) = F + Z/2Z where F is free abelian, and the Z/2Z factor is generated by the middle curve in a Möbius band embedded in S. Orientability of manifolds Let M be a connected topological n-manifold. There are several possible definitions of what it means for M to be orientable. Some of these definitions require that M has extra structure, like being differentiable. Occasionally, must be made into a special case. When more than one of these definitions applies to M, then M is orientable under one definition if and only if it is orientable under the others. Orientability of differentiable manifolds The most intuitive definitions require that M be a differentiable manifold. This means that the transition functions in the atlas of M are C1-functions. Such a function admits a Jacobian determinant. When the Jacobian determinant is positive, the transition function is said to be orientation preserving. An oriented atlas on M is an atlas for which all transition functions are orientation preserving. M is orientable if it admits an oriented atlas. When , an orientation of M is a maximal oriented atlas. (When , an orientation of M is a function .) Orientability and orientations can also be expressed in terms of the tangent bundle. The tangent bundle is a vector bundle, so it is a fiber bundle with structure group . That is, the transition functions of the manifold induce transition functions on the tangent bundle which are fiberwise linear transformations. If the structure group can be reduced to the group of positive determinant matrices, or equivalently if there exists an atlas whose transition functions determine an orientation preserving linear transformation on each tangent space, then the manifold M is orientable. Conversely, M is orientable if and only if the structure group of the tangent bundle can be reduced in this way. Similar observations can be made for the frame bundle. Another way to define orientations on a differentiable manifold is through volume forms. A volume form is a nowhere vanishing section ω of , the top exterior power of the cotangent bundle of M. For example, Rn has a standard volume form given by . Given a volume form on M, the collection of all charts for which the standard volume form pulls back to a positive multiple of ω is an oriented atlas. The existence of a volume form is therefore equivalent to orientability of the manifold. Volume forms and tangent vectors can be combined to give yet another description of orientability. If is a basis of tangent vectors at a point p, then the basis is said to be right-handed if . A transition function is orientation preserving if and only if it sends right-handed bases to right-handed bases. The existence of a volume form implies a reduction of the structure group of the tangent bundle or the frame bundle to . As before, this implies the orientability of M. Conversely, if M is orientable, then local volume forms can be patched together to create a global volume form, orientability being necessary to ensure that the global form is nowhere vanishing. Homology and the orientability of general manifolds At the heart of all the above definitions of orientability of a differentiable manifold is the notion of an orientation preserving transition function. This raises the question of what exactly such transition functions are preserving. They cannot be preserving an orientation of the manifold because an orientation of the manifold is an atlas, and it makes no sense to say that a transition function preserves or does not preserve an atlas of which it is a member. This question can be resolved by defining local orientations. On a one-dimensional manifold, a local orientation around a point p corresponds to a choice of left and right near that point. On a two-dimensional manifold, it corresponds to a choice of clockwise and counter-clockwise. These two situations share the common feature that they are described in terms of top-dimensional behavior near p but not at p. For the general case, let M be a topological n-manifold. A local orientation of M around a point p is a choice of generator of the group To see the geometric significance of this group, choose a chart around p. In that chart there is a neighborhood of p which is an open ball B around the origin O. By the excision theorem, is isomorphic to . The ball B is contractible, so its homology groups vanish except in degree zero, and the space is an -sphere, so its homology groups vanish except in degrees and . A computation with the long exact sequence in relative homology shows that the above homology group is isomorphic to . A choice of generator therefore corresponds to a decision of whether, in the given chart, a sphere around p is positive or negative. A reflection of through the origin acts by negation on , so the geometric significance of the choice of generator is that it distinguishes charts from their reflections. On a topological manifold, a transition function is orientation preserving if, at each point p in its domain, it fixes the generators of . From here, the relevant definitions are the same as in the differentiable case. An oriented atlas is one for which all transition functions are orientation preserving, M is orientable if it admits an oriented atlas, and when , an orientation of M is a maximal oriented atlas. Intuitively, an orientation of M ought to define a unique local orientation of M at each point. This is made precise by noting that any chart in the oriented atlas around p can be used to determine a sphere around p, and this sphere determines a generator of . Moreover, any other chart around p is related to the first chart by an orientation preserving transition function, and this implies that the two charts yield the same generator, whence the generator is unique. Purely homological definitions are also possible. Assuming that M is closed and connected, M is orientable if and only if the nth homology group is isomorphic to the integers Z. An orientation of M is a choice of generator of this group. This generator determines an oriented atlas by fixing a generator of the infinite cyclic group and taking the oriented charts to be those for which pushes forward to the fixed generator. Conversely, an oriented atlas determines such a generator as compatible local orientations can be glued together to give a generator for the homology group . Orientation and cohomology A manifold M is orientable if and only if the first Stiefel–Whitney class vanishes. In particular, if the first cohomology group with Z/2 coefficients is zero, then the manifold is orientable. Moreover, if M is orientable and w1 vanishes, then parametrizes the choices of orientations. This characterization of orientability extends to orientability of general vector bundles over M, not just the tangent bundle. The orientation double cover Around each point of M there are two local orientations. Intuitively, there is a way to move from a local orientation at a point to a local orientation at a nearby point : when the two points lie in the same coordinate chart , that coordinate chart defines compatible local orientations at and . The set of local orientations can therefore be given a topology, and this topology makes it into a manifold. More precisely, let O be the set of all local orientations of M. To topologize O we will specify a subbase for its topology. Let U be an open subset of M chosen such that is isomorphic to Z. Assume that α is a generator of this group. For each p in U, there is a pushforward function . The codomain of this group has two generators, and α maps to one of them. The topology on O is defined so that is open. There is a canonical map that sends a local orientation at p to p. It is clear that every point of M has precisely two preimages under . In fact, is even a local homeomorphism, because the preimages of the open sets U mentioned above are homeomorphic to the disjoint union of two copies of U. If M is orientable, then M itself is one of these open sets, so O is the disjoint union of two copies of M. If M is non-orientable, however, then O is connected and orientable. The manifold O is called the orientation double cover. Manifolds with boundary If M is a manifold with boundary, then an orientation of M is defined to be an orientation of its interior. Such an orientation induces an orientation of ∂M. Indeed, suppose that an orientation of M is fixed. Let be a chart at a boundary point of M which, when restricted to the interior of M, is in the chosen oriented atlas. The restriction of this chart to ∂M is a chart of ∂M. Such charts form an oriented atlas for ∂M. When M is smooth, at each point p of ∂M, the restriction of the tangent bundle of M to ∂M is isomorphic to , where the factor of R is described by the inward pointing normal vector. The orientation of Tp∂M is defined by the condition that a basis of Tp∂M is positively oriented if and only if it, when combined with the inward pointing normal vector, defines a positively oriented basis of TpM. Orientable double cover A closely related notion uses the idea of covering space. For a connected manifold take , the set of pairs where is a point of and is an orientation at ; here we assume is either smooth so we can choose an orientation on the tangent space at a point or we use singular homology to define orientation. Then for every open, oriented subset of we consider the corresponding set of pairs and define that to be an open set of . This gives a topology and the projection sending to is then a 2-to-1 covering map. This covering space is called the orientable double cover, as it is orientable. is connected if and only if is not orientable. Another way to construct this cover is to divide the loops based at a basepoint into either orientation-preserving or orientation-reversing loops. The orientation preserving loops generate a subgroup of the fundamental group which is either the whole group or of index two. In the latter case (which means there is an orientation-reversing path), the subgroup corresponds to a connected double covering; this cover is orientable by construction. In the former case, one can simply take two copies of , each of which corresponds to a different orientation. Orientation of vector bundles A real vector bundle, which a priori has a GL(n) structure group, is called orientable when the structure group may be reduced to , the group of matrices with positive determinant. For the tangent bundle, this reduction is always possible if the underlying base manifold is orientable and in fact this provides a convenient way to define the orientability of a smooth real manifold: a smooth manifold is defined to be orientable if its tangent bundle is orientable (as a vector bundle). Note that as a manifold in its own right, the tangent bundle is always orientable, even over nonorientable manifolds. Related concepts Lorentzian geometry In Lorentzian geometry, there are two kinds of orientability: space orientability and time orientability. These play a role in the causal structure of spacetime. In the context of general relativity, a spacetime manifold is space orientable if, whenever two right-handed observers head off in rocket ships starting at the same spacetime point, and then meet again at another point, they remain right-handed with respect to one another. If a spacetime is time-orientable then the two observers will always agree on the direction of time at both points of their meeting. In fact, a spacetime is time-orientable if and only if any two observers can agree which of the two meetings preceded the other. Formally, the pseudo-orthogonal group O(p,q) has a pair of characters: the space orientation character σ+ and the time orientation character σ−, Their product σ = σ+σ− is the determinant, which gives the orientation character. A space-orientation of a pseudo-Riemannian manifold is identified with a section of the associated bundle where O(M) is the bundle of pseudo-orthogonal frames. Similarly, a time orientation is a section of the associated bundle	Mathematics	Topology	null
187461	https://en.wikipedia.org/wiki/Lenz%27s%20law	Lenz's law	Lenz's law states that the direction of the electric current induced in a conductor by a changing magnetic field is such that the magnetic field created by the induced current opposes changes in the initial magnetic field. It is named after physicist Heinrich Lenz, who formulated it in 1834. It is a qualitative law that specifies the direction of induced current, but states nothing about its magnitude. Lenz's law predicts the direction of many effects in electromagnetism, such as the direction of voltage induced in an inductor or wire loop by a changing current, or the drag force of eddy currents exerted on moving objects in the magnetic field. Lenz's law may be seen as analogous to Newton's third law in classical mechanics and Le Chatelier's principle in chemistry. Definition Lenz's law states that: The current induced in a circuit due to a change in a magnetic field is directed to oppose the change in flux and to exert a mechanical force which opposes the motion. Lenz's law is contained in the rigorous treatment of Faraday's law of induction (the magnitude of EMF induced in a coil is proportional to the rate of change of the magnetic flux), where it finds expression by the negative sign: which indicates that the induced electromotive force and the rate of change in magnetic flux have opposite signs. This means that the direction of the back EMF of an induced field opposes the changing current that is its cause. D.J. Griffiths summarized it as follows: Nature abhors a change in flux. If a change in the magnetic field of current i1 induces another electric current, i2, the direction of i2 is opposite that of the change in i1. If these currents are in two coaxial circular conductors ℓ1 and ℓ2 respectively, and both are initially 0, then the currents i1 and i2 must counter-rotate. The opposing currents will repel each other as a result. Example Magnetic fields from strong magnets can create counter-rotating currents in a copper or aluminium pipe. This is shown by dropping the magnet through the pipe. The descent of the magnet inside the pipe is observably slower than when dropped outside the pipe. When a voltage is generated by a change in magnetic flux according to Faraday's law, the polarity of the induced voltage is such that it produces a current whose magnetic field opposes the change which produces it. The induced magnetic field inside any loop of wire always acts to keep the magnetic flux in the loop constant. The direction of an induced current can be determined using the right-hand rule to show which direction of current flow would create a magnetic field that would oppose the direction of changing flux through the loop. In the examples above, if the flux is increasing, the induced field acts in opposition to it. If it is decreasing, the induced field acts in the direction of the applied field to oppose the change. Detailed interaction of charges in these currents In electromagnetism, when charges move along electric field lines work is done on them, whether it involves storing potential energy (negative work) or increasing kinetic energy (positive work). When net positive work is applied to a charge q1, it gains speed and momentum. The net work on q1 thereby generates a magnetic field whose strength (in units of magnetic flux density (1 tesla = 1 volt-second per square meter)) is proportional to the speed increase of q1. This magnetic field can interact with a neighboring charge q2, passing on this momentum to it, and in return, q1 loses momentum. The charge q2 can also act on q1 in a similar manner, by which it returns some of the momentum that it received from q1. This back-and-forth component of momentum contributes to magnetic inductance. The closer that q1 and q2 are, the greater the effect. When q2 is inside a conductive medium such as a thick slab made of copper or aluminum, it more readily responds to the force applied to it by q1. The energy of q1 is not instantly consumed as heat generated by the current of q2 but is also stored in two opposing magnetic fields. The energy density of magnetic fields tends to vary with the square of the magnetic field's intensity; however, in the case of magnetically non-linear materials such as ferromagnets and superconductors, this relationship breaks down. Conservation of momentum Momentum must be conserved in the process, so if q1 is pushed in one direction, then q2 ought to be pushed in the other direction by the same force at the same time. However, the situation becomes more complicated when the finite speed of electromagnetic wave propagation is introduced (see retarded potential). This means that for a brief period the total momentum of the two charges is not conserved, implying that the difference should be accounted for by momentum in the fields, as asserted by Richard P. Feynman. Famous 19th century electrodynamicist James Clerk Maxwell called this the "electromagnetic momentum". Yet, such a treatment of fields may be necessary when Lenz's law is applied to opposite charges. It is normally assumed that the charges in question have the same sign. If they do not, such as a proton and an electron, the interaction is different. An electron generating a magnetic field would generate an EMF that causes a proton to accelerate in the same direction as the electron. At first, this might seem to violate the law of conservation of momentum, but such an interaction is seen to conserve momentum if the momentum of electromagnetic fields is taken into account.	Physical sciences	Electrodynamics	Physics
187721	https://en.wikipedia.org/wiki/Imprinting%20%28psychology%29	Imprinting (psychology)	In psychology and ethology, imprinting is any kind of phase-sensitive learning (learning occurring at a particular age or a particular life stage) that is rapid and apparently independent of the consequences of behaviour. It was first used to describe situations in which an animal or person learns the characteristics of some stimulus, which is therefore said to be "imprinted" onto the subject. Imprinting is hypothesized to have a critical period. Filial imprinting The best-known form of imprinting is filial imprinting, in which a young animal narrows its social preferences to an object (typically a parent) as a result of exposure to that object. It is most obvious in nidifugous birds, which imprint on their parents and then follow them around. It was first reported in domestic chickens, by Sir Thomas More in 1516 as described in his treatise Utopia, 350years earlier than by the 19th-century amateur biologist Douglas Spalding. It was rediscovered by the early ethologist Oskar Heinroth, and studied extensively and popularized by his disciple Konrad Lorenz working with greylag geese. Lorenz demonstrated how incubator-hatched geese would imprint on the first suitable moving stimulus they saw within what he called a "critical period" between 13 and 16 hours shortly after hatching. For example, the goslings would imprint on Lorenz himself (to be more specific, on his wading boots), and he is often depicted being followed by a gaggle of geese who had imprinted on him. Lorenz also found that the geese could imprint on inanimate objects. In one notable experiment, they followed a box placed on a model train in circles around the track. Filial imprinting is not restricted to non-human animals that are able to follow their parents, however. The filial imprinting of birds was a primary technique used to create the movie Winged Migration (Le Peuple Migrateur), which contains a great deal of footage of migratory birds in flight. The birds imprinted on handlers, who wore yellow jackets and honked horns constantly. The birds were then trained to fly along with a variety of aircraft, primarily ultralights. The Italian hang-glider pilot Angelo d'Arrigo extended this technique. D'Arrigo noted that the flight of a non-motorised hang-glider is very similar to the flight patterns of migratory birds; both use updrafts of hot air (thermal currents) to gain altitude that then permits soaring flight over distance. He used this to reintroduce threatened species of raptors. Because birds hatched in captivity have no mentor birds to teach them traditional migratory routes, D'Arrigo hatched chicks under the wing of his glider and they imprinted on him. Then, he taught the fledglings to fly and to hunt. The young birds followed him not only on the ground (as with Lorenz) but also in the air as he took the path of various migratory routes. He flew across the Sahara and over the Mediterranean Sea to Sicily with eagles, from Siberia to Iran (5,500 km) with a flock of Siberian cranes, and over Mount Everest with Nepalese eagles. In 2006, he worked with a condor in South America. In a similar project, orphaned Canada geese were trained to their normal migration route by the Canadian ultralight enthusiast Bill Lishman, as shown in the fact-based movie drama Fly Away Home. Chicks of domestic chickens prefer to be near large groups of objects that they have imprinted on. This behaviour was used to determine that very young chicks of a few days old have rudimentary counting skills. In a series of experiments, they were made to imprint on plastic balls and could figure out which of two groups of balls hidden behind screens had the most balls. American coot mothers have the ability to recognize their chicks by imprinting on cues from the first chick that hatches. This allows mothers to distinguish their chicks from parasitic chicks. The peregrine falcon has also been known to imprint on specific structures for their breeding grounds such as cliff sides and bridges and thus will favour that location for breeding. Sexual imprinting Sexual imprinting is the process by which a young animal learns the characteristics of a desirable mate. For example, male zebra finches appear to prefer mates with the appearance of the female bird that rears them. Sexual attraction to humans can develop in non-human mammals or birds as a result of sexual imprinting when reared from young by humans. One example is London Zoo female giant panda Chi Chi. When taken to Moscow Zoo for mating with the male giant panda An An, she refused his attempts to mate with her, but made a full sexual self-presentation to a zookeeper. It commonly occurs in falconry birds reared from hatching by humans. Such birds are called "imprints" in falconry. When an imprint must be bred from, the breeder lets the male bird copulate with their head while they are wearing a special hat with pockets to catch the male bird's semen. The breeder then courts a suitable imprint female bird (including offering food, if it is part of that species's normal courtship). At "copulation", the breeder puts the flat of one hand on the female bird's back to represent the weight of a male bird, and with the other hand uses a pipette, or a hypodermic syringe without a needle, to squirt the semen into the female's cloaca. Sexual imprinting on inanimate objects is a popular theory concerning the development of sexual fetishism. For example, according to this theory, imprinting on shoes or boots (as with Konrad Lorenz's geese) would be the cause of shoe fetishism. Limbic imprinting Some suggest that prenatal, perinatal and post-natal experiences leave imprints upon the limbic system, causing lifelong effects and this process is identified as limbic imprinting. The term is also described as the human emotional map, deep-seated beliefs, and values that are stored in the brain's limbic system and govern people's lives at the subconscious level. It is one of the suggested explanations for the claim that the experiences of an infant, particularly during the first two years of life, contribute to a person's lifelong psychological development. Imprinted genes can have astounding effects on body size, brain size, and the process in which the brain organizes its processes. Evolutionary trends within the animal kingdom have been shown to show substantive increase in the forebrain particularly towards the limbic system. This evolution has even been thought of to have a mutative effect on the brain size trickling down the human ancestry. Westermarck effect Reverse sexual imprinting is also seen in instances where two people who live in domestic proximity during the first few years in the life of either one become desensitized to later close sexual attraction to each other. This phenomenon, known as the Westermarck effect, was first formally described by Finnish anthropologist Edvard Westermarck in his book The History of Human Marriage (1891). The Westermarck effect has since been observed in many places and cultures, including in the Israeli kibbutz system, and the Chinese shim-pua marriage customs, as well as in biological-related families. In the case of the Israeli kibbutzim (collective farms), children were reared somewhat communally in peer groups, based on age, not biological relation. A study of the marriage patterns of these children later in life revealed that out of the nearly 3,000 marriages that occurred across the kibbutz system, only fourteen were between children from the same peer group. Of those fourteen, none had been reared together during the first six years of life. This result provides evidence not only that the Westermarck effect is demonstrable but that it operates during the period from birth to the age of six. However, Eran Shor and Dalit Simchai claimed that the case of the kibbutzim actually provides little support for the Westermarck effect. When proximity during this critical period does not occur—for example, where a brother and sister are brought up separately, never meeting one another—they may find one another highly sexually attractive when they meet as adults. This phenomenon is known as genetic sexual attraction. This observation supports the hypothesis that the Westermarck effect evolved because it suppressed inbreeding. This attraction may also be seen with cousin couples. Sigmund Freud argued that as children, members of the same family naturally lust for one another, making it necessary for societies to create incest taboos, but Westermarck argued the reverse, that the taboos themselves arise naturally as products of innate attitudes. Steven Pinker has written that Freud's conception of an urge to incest may have derived from Freud's own erotic reaction to his mother as a boy (attested in Freud's own writings), and speculates that Freud's reaction may have been due to lack of intimacy with his mother in early childhood, as Freud was wet-nursed. Baby duck syndrome In human–computer interaction, baby duck syndrome denotes the tendency for computer users to "imprint" on the first system they learn, then judge other systems by their similarity to that first system. The result is that "users generally prefer systems similar to those they learned on and dislike unfamiliar systems". The issue may present itself relatively early in a computer user's experience, and it has been observed to impede education of students in new software systems or user interfaces.	Biology and health sciences	Ethology	Biology
187769	https://en.wikipedia.org/wiki/Pseudoscorpion	Pseudoscorpion	Pseudoscorpions, also known as false scorpions or book scorpions, are small, scorpion-like arachnids belonging to the order Pseudoscorpiones, also known as Pseudoscorpionida or Chelonethida. Pseudoscorpions are generally beneficial to humans because they prey on clothes moth larvae, carpet beetle larvae, booklice, ants, mites, and small flies. They are common in many environments, but they are rarely noticed due to their small size. When people see pseudoscorpions, especially indoors, they often mistake them for ticks or small spiders. Pseudoscorpions often carry out phoresis, a form of commensalism in which one organism uses another for the purpose of transport. Characteristics Pseudoscorpions belong to the class Arachnida. They are small arachnids with a flat, pear-shaped body, and pincer-like pedipalps that resemble those of scorpions. They usually range from in length. The largest known species is Garypus titanius of Ascension Island at up to . Range is generally smaller at an average of . A pseudoscorpion has eight legs with five to seven segments each; the number of fused segments is used to distinguish families and genera. They have two very long pedipalps with palpal chelae (pincers), which strongly resemble the pincers found on a scorpion. The pedipalps generally consist of an immobile "hand" and mobile "finger", the latter controlled by an adductor muscle. Members of the clade Iocheirata, which contains the majority of pseudoscorpions, are venomous, with a venom gland and duct usually located in the mobile finger; the venom is used to immobilize the pseudoscorpion's prey. During digestion, pseudoscorpions exude a mildly corrosive fluid over the prey, then ingest the liquefied remains. The abdomen, referred to as the opisthosoma, is made up of twelve segments, each protected by sclerotized plates (called tergites above and sternites below). The abdomen is short and rounded at the rear, rather than extending into a segmented tail and stinger like true scorpions. The color of the body can be yellowish-tan to dark-brown, with the paired claws often a contrasting color. They may have two, four or no eyes. Pseudoscorpions spin silk from a gland in their jaws to make disk-shaped cocoons for mating, molting, or waiting out cold weather, but they do not have book lungs like true scorpions and the Tetrapulmonata. Instead, they breathe exclusively through tracheae, which open laterally through two pairs of spiracles on the posterior margins of the sternites of abdominal segments 3 and 4. Behavior The male produces a spermatophore which is attached to the substrate and is picked up by the female. Members of the Cheliferoidea (Atemnidae, Cheliferidae, Chernetidae and Withiidae) have an elaborate mating dance, which ends with the male navigating the female over his spermatophore. In Cheliferidae, the male also uses his forelegs to open the female genital operculum, and after she has mounted the packet of sperm, assisting the spermatophore's entry by pushing it into her genital opening. Females in species that possess a spermatheca (sperm storing organ) can store the sperm for a longer period of time before fertilizing the eggs, but species without the organ fertilize the eggs shortly after mating. The female carries the fertilized eggs in a brood pouch attached to her abdomen. Between 2 and 50 young are hatched in a single brood, with more than one brood per year possible. The young go through three molts called the protonymph, deutonymph and tritonymph. The developing embryo and the protonymph, which remain attached to the mother, is nourished by a ‘milk’ produced by her ovary. Many species molt in a small, silken igloo that protects them from enemies during this vulnerable period. After reaching adulthood they no longer molt, and will live for 2–3 years. They are active in the warm months of the year, overwintering in silken cocoons when the weather grows cold. Smaller species live in debris and humus. Some species are arboreal, while others are phagophiles, eating parasites in an example of cleaning symbiosis. Some species are phoretic, others may sometimes be found feeding on mites under the wing covers of certain beetles. Distribution More than 3,300 species of pseudoscorpions are recorded in more than 430 genera, with more being discovered on a regular basis. They range worldwide, even in temperate to cold regions such as Northern Ontario and above the timberline in Wyoming's Rocky Mountains in the United States and the Jenolan Caves of Australia, but have their most dense and diverse populations in the tropics and subtropics, where they spread even to island territories such as the Canary Islands, where around 25 endemic species have been found. There are also two endemic species on the Maltese Islands. Species have been found under tree bark, in leaf and pine litter, in soil, in tree hollows, under stones, in caves such as the Movile Cave, at the seashore in the intertidal zone, and within fractured rocks. Chelifer cancroides is the species most commonly found in homes, where it is often observed in rooms with dusty books. There, the tiny animals () can find their food such as booklice and house dust mites. They enter homes by riding insects (phoresy) larger than themselves, or are brought in with firewood. Evolution The oldest known fossil pseudoscorpion, Dracochela deprehendor is known from cuticle fragments of nymphs found in the Panther Mountain Formation near Gilboa in New York, dating to the mid-Devonian, around 383 million years ago. It has all of the traits of a modern pseudoscorpion, indicating that the order evolved very early in the history of land animals. Its morphology suggests that it is more primitive than any living pseudoscorpion. As with most other arachnid orders, the pseudoscorpions have changed very little since they first appeared, retaining almost all the features of their original form. After the Devonian fossils, almost no other fossils of pseudoscorpions are known for over 250 million years until Cretaceous fossils in amber, all belonging to modern families, suggesting that the major diversification of pseudoscorpions had already taken place by this time. The only fossil from this time gap is Archaeofeaella from the Triassic of Ukraine, approximately 227 million years ago, which is suggested to be an early relative of the family Feaellidae. Historical references Pseudoscorpions were first described by Aristotle, who probably found them among scrolls in a library where they would have been feeding on booklice. Robert Hooke referred to a "Land-Crab" in his 1665 work Micrographia. Another reference in the 1780s, when George Adams wrote of "a lobster-insect, spied by some labouring men who were drinking their porter, and borne away by an ingenious gentleman, who brought it to my lodging." Classification The following taxon numbers are calculated as of the end of 2023. Atemnidae Kishida, 1929 (21 genera, 194 species) Bochicidae Chamberlin, 1930 (12 genera, 44 species) Cheiridiidae Hansen, 1894 (9 genera, 81 species) Cheliferidae Risso, 1827 (64 genera, 312 species) Chernetidae Menge, 1855 (120 genera, 728 species) Chthoniidae Daday, 1888 (54 genera, 909 species) Feaellidae Ellingsen, 1906 (8 genus, 37 species) Garypidae Simon, 1879 (11 genera, 110 species) Garypinidae Daday, 1888 (21 genera, 94 species) Geogarypidae Chamberlin, 1930 (2 genera, 81 species) Gymnobisiidae Beier, 1947 (4 genera, 17 species) Hyidae Chamberlin, 1930 (2 genera, 41 species) Ideoroncidae Chamberlin, 1930 (15 genera, 86 species) Larcidae Harvey, 1992 (1 genus, 15 species) Menthidae Chamberlin, 1930 (5 genera, 12 species) Neobisiidae Chamberlin, 1930 (34 genera, 748 species) Olpiidae Banks, 1895 (24 genera, 211 species) Parahyidae Harvey, 1992 (1 genus, 1 species) Pseudochiridiidae Chamberlin, 1923 (2 genera, 13 species) Pseudogarypidae Chamberlin, 1923 (2 genera, 12 species) Pseudotyrannochthoniidae Beier, 1932 (6 genera, 80 species) Sternophoridae Chamberlin, 1923 (3 genera, 21 species) Syarinidae Chamberlin, 1930 (18 genera, 125 species) Withiidae Chamberlin, 1931 (37 genera, 170 species) †Dracochelidae Schawaller, Shear & Bonamo, 1991 (1 genus, 1 species) Cladogram After Benavides et al., 2019, with historic taxonomic groups from Harvey (1992).	Biology and health sciences	Arachnids	null
187813	https://en.wikipedia.org/wiki/Transport%20Layer%20Security	Transport Layer Security	Transport Layer Security (TLS) is a cryptographic protocol designed to provide communications security over a computer network, such as the Internet. The protocol is widely used in applications such as email, instant messaging, and voice over IP, but its use in securing HTTPS remains the most publicly visible. The TLS protocol aims primarily to provide security, including privacy (confidentiality), integrity, and authenticity through the use of cryptography, such as the use of certificates, between two or more communicating computer applications. It runs in the presentation layer and is itself composed of two layers: the TLS record and the TLS handshake protocols. The closely related Datagram Transport Layer Security (DTLS) is a communications protocol that provides security to datagram-based applications. In technical writing, references to "(D)TLS" are often seen when it applies to both versions. TLS is a proposed Internet Engineering Task Force (IETF) standard, first defined in 1999, and the current version is TLS 1.3, defined in August 2018. TLS builds on the now-deprecated SSL (Secure Sockets Layer) specifications (1994, 1995, 1996) developed by Netscape Communications for adding the HTTPS protocol to their Netscape Navigator web browser. Description Client-server applications use the TLS protocol to communicate across a network in a way designed to prevent eavesdropping and tampering. Since applications can communicate either with or without TLS (or SSL), it is necessary for the client to request that the server set up a TLS connection. One of the main ways of achieving this is to use a different port number for TLS connections. Port 80 is typically used for unencrypted HTTP traffic while port 443 is the common port used for encrypted HTTPS traffic. Another mechanism is to make a protocol-specific STARTTLS request to the server to switch the connection to TLS – for example, when using the mail and news protocols. Once the client and server have agreed to use TLS, they negotiate a stateful connection by using a handshaking procedure (see ). The protocols use a handshake with an asymmetric cipher to establish not only cipher settings but also a session-specific shared key with which further communication is encrypted using a symmetric cipher. During this handshake, the client and server agree on various parameters used to establish the connection's security: The handshake begins when a client connects to a TLS-enabled server requesting a secure connection and the client presents a list of supported cipher suites (ciphers and hash functions). From this list, the server picks a cipher and hash function that it also supports and notifies the client of the decision. The server usually then provides identification in the form of a digital certificate. The certificate contains the server name, the trusted certificate authority (CA) that vouches for the authenticity of the certificate, and the server's public encryption key. The client confirms the validity of the certificate before proceeding. To generate the session keys used for the secure connection, the client either: encrypts a random number (PreMasterSecret) with the server's public key and sends the result to the server (which only the server should be able to decrypt with its private key); both parties then use the random number to generate a unique session key for subsequent encryption and decryption of data during the session, or uses Diffie–Hellman key exchange (or its variant elliptic-curve DH) to securely generate a random and unique session key for encryption and decryption that has the additional property of forward secrecy: if the server's private key is disclosed in future, it cannot be used to decrypt the current session, even if the session is intercepted and recorded by a third party. This concludes the handshake and begins the secured connection, which is encrypted and decrypted with the session key until the connection closes. If any one of the above steps fails, then the TLS handshake fails and the connection is not created. TLS and SSL do not fit neatly into any single layer of the OSI model or the TCP/IP model. TLS runs "on top of some reliable transport protocol (e.g., TCP)," which would imply that it is above the transport layer. It serves encryption to higher layers, which is normally the function of the presentation layer. However, applications generally use TLS as if it were a transport layer, even though applications using TLS must actively control initiating TLS handshakes and handling of exchanged authentication certificates. When secured by TLS, connections between a client (e.g., a web browser) and a server (e.g., wikipedia.org) will have all of the following properties: The connection is private (or has confidentiality) because a symmetric-key algorithm is used to encrypt the data transmitted. The keys for this symmetric encryption are generated uniquely for each connection and are based on a shared secret that was negotiated at the start of the session. The server and client negotiate the details of which encryption algorithm and cryptographic keys to use before the first byte of data is transmitted (see below). The negotiation of a shared secret is both secure (the negotiated secret is unavailable to eavesdroppers and cannot be obtained, even by an attacker who places themselves in the middle of the connection) and reliable (no attacker can modify the communications during the negotiation without being detected). The identity of the communicating parties can be authenticated using public-key cryptography. This authentication is required for the server and optional for the client. The connection is reliable (or has integrity) because each message transmitted includes a message integrity check using a message authentication code to prevent undetected loss or alteration of the data during transmission. TLS supports many different methods for exchanging keys, encrypting data, and authenticating message integrity. As a result, secure configuration of TLS involves many configurable parameters, and not all choices provide all of the privacy-related properties described in the list above (see the tables below § Key exchange, § Cipher security, and ). Attempts have been made to subvert aspects of the communications security that TLS seeks to provide, and the protocol has been revised several times to address these security threats. Developers of web browsers have repeatedly revised their products to defend against potential security weaknesses after these were discovered (see TLS/SSL support history of web browsers). Datagram Transport Layer Security Datagram Transport Layer Security, abbreviated DTLS, is a related communications protocol providing security to datagram-based applications by allowing them to communicate in a way designed to prevent eavesdropping, tampering, or message forgery. The DTLS protocol is based on the stream-oriented Transport Layer Security (TLS) protocol and is intended to provide similar security guarantees. However, unlike TLS, it can be used with most datagram oriented protocols including User Datagram Protocol (UDP), Datagram Congestion Control Protocol (DCCP), Control And Provisioning of Wireless Access Points (CAPWAP), Stream Control Transmission Protocol (SCTP) encapsulation, and Secure Real-time Transport Protocol (SRTP). As the DTLS protocol datagram preserves the semantics of the underlying transport—the application it does not suffer from the delays associated with stream protocols, however the application has to deal with packet reordering, loss of datagram and data larger than the size of a datagram network packet. Because DTLS uses UDP or SCTP rather than TCP, it avoids the TCP meltdown problem, when being used to create a VPN tunnel. The original 2006 release of DTLS version 1.0 was not a standalone document. It was given as a series of deltas to TLS 1.1. Similarly the follow-up 2012 release of DTLS is a delta to TLS 1.2. It was given the version number of DTLS 1.2 to match its TLS version. Lastly, the 2022 DTLS 1.3 is a delta to TLS 1.3. Like the two previous versions, DTLS 1.3 is intended to provide "equivalent security guarantees [to TLS 1.3] with the exception of order protection/non-replayability". Many VPN clients including Cisco AnyConnect & InterCloud Fabric, OpenConnect, ZScaler tunnel, F5 Networks Edge VPN Client, and Citrix Systems NetScaler use DTLS to secure UDP traffic. In addition all modern web browsers support DTLS-SRTP for WebRTC. History and development Secure Data Network System The Transport Layer Security Protocol (TLS), together with several other basic network security platforms, was developed through a joint initiative begun in August 1986, among the National Security Agency, the National Bureau of Standards, the Defense Communications Agency, and twelve communications and computer corporations who initiated a special project called the Secure Data Network System (SDNS). The program was described in September 1987 at the 10th National Computer Security Conference in an extensive set of published papers. The innovative research program focused on designing the next generation of secure computer communications network and product specifications to be implemented for applications on public and private internets. It was intended to complement the rapidly emerging new OSI internet standards moving forward both in the U.S. government's GOSIP Profiles and in the huge ITU-ISO JTC1 internet effort internationally. Originally known as the SP4 protocol, it was renamed TLS and subsequently published in 1995 as international standard ITU-T X.274\|ISO/IEC 10736:1995. Secure Network Programming (SNP) Early research efforts towards transport layer security included the Secure Network Programming (SNP) application programming interface (API), which in 1993 explored the approach of having a secure transport layer API closely resembling Berkeley sockets, to facilitate retrofitting pre-existing network applications with security measures. SNP was published and presented in the 1994 USENIX Summer Technical Conference. The SNP project was funded by a grant from NSA to Professor Simon Lam at UT-Austin in 1991. Secure Network Programming won the 2004 ACM Software System Award. Simon Lam was inducted into the Internet Hall of Fame for "inventing secure sockets and implementing the first secure sockets layer, named SNP, in 1993." SSL 1.0, 2.0, and 3.0 Netscape developed the original SSL protocols, and Taher Elgamal, chief scientist at Netscape Communications from 1995 to 1998, has been described as the "father of SSL". SSL version 1.0 was never publicly released because of serious security flaws in the protocol. Version 2.0, after being released in February 1995 was quickly found to contain a number of security and usability flaws. It used the same cryptographic keys for message authentication and encryption. It had a weak MAC construction that used the MD5 hash function with a secret prefix, making it vulnerable to length extension attacks. It also provided no protection for either the opening handshake or an explicit message close, both of which meant man-in-the-middle attacks could go undetected. Moreover, SSL 2.0 assumed a single service and a fixed domain certificate, conflicting with the widely used feature of virtual hosting in Web servers, so most websites were effectively impaired from using SSL. These flaws necessitated the complete redesign of the protocol to SSL version 3.0. Released in 1996, it was produced by Paul Kocher working with Netscape engineers Phil Karlton and Alan Freier, with a reference implementation by Christopher Allen and Tim Dierks of Certicom. Newer versions of SSL/TLS are based on SSL 3.0. The 1996 draft of SSL 3.0 was published by IETF as a historical document in . SSL 2.0 was deprecated in 2011 by . In 2014, SSL 3.0 was found to be vulnerable to the POODLE attack that affects all block ciphers in SSL; RC4, the only non-block cipher supported by SSL 3.0, is also feasibly broken as used in SSL 3.0. SSL 3.0 was deprecated in June 2015 by . TLS 1.0 TLS 1.0 was first defined in in January 1999 as an upgrade of SSL Version 3.0, and written by Christopher Allen and Tim Dierks of Certicom. As stated in the RFC, "the differences between this protocol and SSL 3.0 are not dramatic, but they are significant enough to preclude interoperability between TLS 1.0 and SSL 3.0". Tim Dierks later wrote that these changes, and the renaming from "SSL" to "TLS", were a face-saving gesture to Microsoft, "so it wouldn't look [like] the IETF was just rubberstamping Netscape's protocol". The PCI Council suggested that organizations migrate from TLS 1.0 to TLS 1.1 or higher before June 30, 2018. In October 2018, Apple, Google, Microsoft, and Mozilla jointly announced they would deprecate TLS 1.0 and 1.1 in March 2020. TLS 1.0 and 1.1 were formally deprecated in in March 2021. TLS 1.1 TLS 1.1 was defined in RFC 4346 in April 2006. It is an update from TLS version 1.0. Significant differences in this version include: Added protection against cipher-block chaining (CBC) attacks. The implicit initialization vector (IV) was replaced with an explicit IV. Change in handling of padding errors. Support for IANA registration of parameters. Support for TLS versions 1.0 and 1.1 was widely deprecated by web sites around 2020, disabling access to Firefox versions before 24 and Chromium-based browsers before 29. TLS 1.2 TLS 1.2 was defined in in August 2008. It is based on the earlier TLS 1.1 specification. Major differences include: The MD5 and SHA-1 combination in the pseudorandom function (PRF) was replaced with SHA-256, with an option to use cipher suite specified PRFs. The MD5 and SHA-1 combination in the finished message hash was replaced with SHA-256, with an option to use cipher suite specific hash algorithms. However, the size of the hash in the finished message must still be at least 96 bits. The MD5 and SHA-1 combination in the digitally signed element was replaced with a single hash negotiated during handshake, which defaults to SHA-1. Enhancement in the client's and server's ability to specify which hashes and signature algorithms they accept. Expansion of support for authenticated encryption ciphers, used mainly for Galois/Counter Mode (GCM) and CCM mode of Advanced Encryption Standard (AES) encryption. TLS Extensions definition and AES cipher suites were added. All TLS versions were further refined in in March 2011, removing their backward compatibility with SSL such that TLS sessions never negotiate the use of Secure Sockets Layer (SSL) version 2.0. There is currently no formal date for TLS 1.2 to be deprecated. The specifications for TLS 1.2 became redefined as well by the Standards Track Document to keep it as secure as possible; it is to be seen as a failover protocol now, meant only to be negotiated with clients which are unable to talk over TLS 1.3 (The original RFC 5246 definition for TLS 1.2 is since then obsolete). TLS 1.3 TLS 1.3 was defined in RFC 8446 in August 2018. It is based on the earlier TLS 1.2 specification. Major differences from TLS 1.2 include: Separating key agreement and authentication algorithms from the cipher suites Removing support for weak and less-used named elliptic curves Removing support for MD5 and SHA-224 cryptographic hash functions Requiring digital signatures even when a previous configuration is used Integrating HKDF and the semi-ephemeral DH proposal Replacing resumption with PSK and tickets Supporting 1-RTT handshakes and initial support for 0-RTT Mandating perfect forward secrecy, by means of using ephemeral keys during the (EC)DH key agreement Dropping support for many insecure or obsolete features including compression, renegotiation, non-AEAD ciphers, null ciphers, non-PFS key exchange (among which are static RSA and static DH key exchanges), custom DHE groups, EC point format negotiation, Change Cipher Spec protocol, Hello message UNIX time, and the length field AD input to AEAD ciphers Prohibiting SSL or RC4 negotiation for backwards compatibility Integrating use of session hash Deprecating use of the record layer version number and freezing the number for improved backwards compatibility Moving some security-related algorithm details from an appendix to the specification and relegating ClientKeyShare to an appendix Adding the ChaCha20 stream cipher with the Poly1305 message authentication code Adding the Ed25519 and Ed448 digital signature algorithms Adding the x25519 and x448 key exchange protocols Adding support for sending multiple OCSP responses Encrypting all handshake messages after the ServerHello Network Security Services (NSS), the cryptography library developed by Mozilla and used by its web browser Firefox, enabled TLS 1.3 by default in February 2017. TLS 1.3 support was subsequently added — but due to compatibility issues for a small number of users, not automatically enabled — to Firefox 52.0, which was released in March 2017. TLS 1.3 was enabled by default in May 2018 with the release of Firefox 60.0. Google Chrome set TLS 1.3 as the default version for a short time in 2017. It then removed it as the default, due to incompatible middleboxes such as Blue Coat web proxies. The intolerance of the new version of TLS was protocol ossification; middleboxes had ossified the protocol's version parameter. As a result, version 1.3 mimics the wire image of version 1.2. This change occurred very late in the design process, only having been discovered during browser deployment. The discovery of this intolerance also led to the prior version negotiation strategy, where the highest matching version was picked, being abandoned due to unworkable levels of ossification. 'Greasing' an extension point, where one protocol participant claims support for non-existent extensions to ensure that unrecognised-but-actually-existent extensions are tolerated and so to resist ossification, was originally designed for TLS, but it has since been adopted elsewhere. During the IETF 100 Hackathon, which took place in Singapore in 2017, the TLS Group worked on adapting open-source applications to use TLS 1.3. The TLS group was made up of individuals from Japan, United Kingdom, and Mauritius via the cyberstorm.mu team. This work was continued in the IETF 101 Hackathon in London, and the IETF 102 Hackathon in Montreal. wolfSSL enabled the use of TLS 1.3 as of version 3.11.1, released in May 2017. As the first commercial TLS 1.3 implementation, wolfSSL 3.11.1 supported Draft 18 and now supports Draft 28, the final version, as well as many older versions. A series of blogs were published on the performance difference between TLS 1.2 and 1.3. In September 2018, the popular OpenSSL project released version 1.1.1 of its library, in which support for TLS 1.3 was "the headline new feature". Support for TLS 1.3 was added to Secure Channel (schannel) for the releases of Windows 11 and Windows Server 2022. Enterprise Transport Security The Electronic Frontier Foundation praised TLS 1.3 and expressed concern about the variant protocol Enterprise Transport Security (ETS) that intentionally disables important security measures in TLS 1.3. Originally called Enterprise TLS (eTLS), ETS is a published standard known as the 'ETSI TS103523-3', "Middlebox Security Protocol, Part3: Enterprise Transport Security". It is intended for use entirely within proprietary networks such as banking systems. ETS does not support forward secrecy so as to allow third-party organizations connected to the proprietary networks to be able to use their private key to monitor network traffic for the detection of malware and to make it easier to conduct audits. Despite the claimed benefits, the EFF warned that the loss of forward secrecy could make it easier for data to be exposed along with saying that there are better ways to analyze traffic. Digital certificates A digital certificate certifies the ownership of a public key by the named subject of the certificate, and indicates certain expected usages of that key. This allows others (relying parties) to rely upon signatures or on assertions made by the private key that corresponds to the certified public key. Keystores and trust stores can be in various formats, such as .pem, .crt, .pfx, and .jks. Certificate authorities TLS typically relies on a set of trusted third-party certificate authorities to establish the authenticity of certificates. Trust is usually anchored in a list of certificates distributed with user agent software, and can be modified by the relying party. According to Netcraft, who monitors active TLS certificates, the market-leading certificate authority (CA) has been Symantec since the beginning of their survey (or VeriSign before the authentication services business unit was purchased by Symantec). As of 2015, Symantec accounted for just under a third of all certificates and 44% of the valid certificates used by the 1 million busiest websites, as counted by Netcraft. In 2017, Symantec sold its TLS/SSL business to DigiCert. In an updated report, it was shown that IdenTrust, DigiCert, and Sectigo are the top 3 certificate authorities in terms of market share since May 2019. As a consequence of choosing X.509 certificates, certificate authorities and a public key infrastructure are necessary to verify the relation between a certificate and its owner, as well as to generate, sign, and administer the validity of certificates. While this can be more convenient than verifying the identities via a web of trust, the 2013 mass surveillance disclosures made it more widely known that certificate authorities are a weak point from a security standpoint, allowing man-in-the-middle attacks (MITM) if the certificate authority cooperates (or is compromised). Importance of SSL Certificates Encryption: SSL certificates encrypt data sent between a web server and a user’s browser, ensuring that sensitive information is protected throughout transmission. This encryption technology stops unauthorized parties from intercepting and interpreting data, so protecting it from possible risks such as hacking or data breaches. Authentication: SSL certificates also offer authentication, certifying the integrity of a website and that visitors are connecting to the correct server rather than a malicious impostor. This authentication method helps consumers gain trust by ensuring that they are dealing with a trustworthy and secure website. Integrity: Another important role of SSL certificates is to ensure data integrity. SSL uses cryptographic techniques to verify that data communicated between the server and the browser is intact and unmodified during transit. This keeps malevolent actors from interfering with the data, ensuring its integrity and trustworthiness. Algorithms Key exchange or key agreement Before a client and server can begin to exchange information protected by TLS, they must securely exchange or agree upon an encryption key and a cipher to use when encrypting data (see ). Among the methods used for key exchange/agreement are: public and private keys generated with RSA (denoted TLS_RSA in the TLS handshake protocol), Diffie–Hellman (TLS_DH), ephemeral Diffie–Hellman (TLS_DHE), elliptic-curve Diffie–Hellman (TLS_ECDH), ephemeral elliptic-curve Diffie–Hellman (TLS_ECDHE), anonymous Diffie–Hellman (TLS_DH_anon), pre-shared key (TLS_PSK) and Secure Remote Password (TLS_SRP). The TLS_DH_anon and TLS_ECDH_anon key agreement methods do not authenticate the server or the user and hence are rarely used because those are vulnerable to man-in-the-middle attacks. Only TLS_DHE and TLS_ECDHE provide forward secrecy. Public key certificates used during exchange/agreement also vary in the size of the public/private encryption keys used during the exchange and hence the robustness of the security provided. In July 2013, Google announced that it would no longer use 1024-bit public keys and would switch instead to 2048-bit keys to increase the security of the TLS encryption it provides to its users because the encryption strength is directly related to the key size. Cipher	Technology	Computer security	null
187849	https://en.wikipedia.org/wiki/Fitness%20%28biology%29	Fitness (biology)	Fitness (often denoted or ω in population genetics models) is a quantitative representation of individual reproductive success. It is also equal to the average contribution to the gene pool of the next generation, made by the same individuals of the specified genotype or phenotype. Fitness can be defined either with respect to a genotype or to a phenotype in a given environment or time. The fitness of a genotype is manifested through its phenotype, which is also affected by the developmental environment. The fitness of a given phenotype can also be different in different selective environments. With asexual reproduction, it is sufficient to assign fitnesses to genotypes. With sexual reproduction, recombination scrambles alleles into different genotypes every generation; in this case, fitness values can be assigned to alleles by averaging over possible genetic backgrounds. Natural selection tends to make alleles with higher fitness more common over time, resulting in Darwinian evolution. The term "Darwinian fitness" can be used to make clear the distinction with physical fitness. Fitness does not include a measure of survival or life-span; Herbert Spencer's well-known phrase "survival of the fittest" should be interpreted as: "Survival of the form (phenotypic or genotypic) that will leave the most copies of itself in successive generations." Inclusive fitness differs from individual fitness by including the ability of an allele in one individual to promote the survival and/or reproduction of other individuals that share that allele, in preference to individuals with a different allele. To avoid double counting, inclusive fitness excludes the contribution of other individuals to the survival and reproduction of the focal individual. One mechanism of inclusive fitness is kin selection. Fitness as propensity Fitness is often defined as a propensity or probability, rather than the actual number of offspring. For example, according to Maynard Smith, "Fitness is a property, not of an individual, but of a class of individuals—for example homozygous for allele A at a particular locus. Thus the phrase 'expected number of offspring' means the average number, not the number produced by some one individual. If the first human infant with a gene for levitation were struck by lightning in its pram, this would not prove the new genotype to have low fitness, but only that the particular child was unlucky." Alternatively, "the fitness of the individual—having an array x of phenotypes—is the probability, s(x), that the individual will be included among the group selected as parents of the next generation." Models of fitness In order to avoid the complications of sex and recombination, the concept of fitness is presented below in the restricted setting of an asexual population without genetic recombination. Thus, fitnesses can be assigned directly to genotypes. There are two commonly used operationalizations of fitness – absolute fitness and relative fitness. Absolute fitness The absolute fitness () of a genotype is defined as the proportional change in the abundance of that genotype over one generation attributable to selection. For example, if is the abundance of a genotype in generation in an infinitely large population (so that there is no genetic drift), and neglecting the change in genotype abundances due to mutations, then . An absolute fitness larger than 1 indicates growth in that genotype's abundance; an absolute fitness smaller than 1 indicates decline. Relative fitness Whereas absolute fitness determines changes in genotype abundance, relative fitness () determines changes in genotype frequency. If is the total population size in generation , and the relevant genotype's frequency is , then , where is the mean relative fitness in the population (again setting aside changes in frequency due to drift and mutation). Relative fitnesses only indicate the change in prevalence of different genotypes relative to each other, and so only their values relative to each other are important; relative fitnesses can be any nonnegative number, including 0. It is often convenient to choose one genotype as a reference and set its relative fitness to 1. Relative fitness is used in the standard Wright–Fisher and Moran models of population genetics. Absolute fitnesses can be used to calculate relative fitness, since (we have used the fact that , where is the mean absolute fitness in the population). This implies that , or in other words, relative fitness is proportional to . It is not possible to calculate absolute fitnesses from relative fitnesses alone, since relative fitnesses contain no information about changes in overall population abundance . Assigning relative fitness values to genotypes is mathematically appropriate when two conditions are met: first, the population is at demographic equilibrium, and second, individuals vary in their birth rate, contest ability, or death rate, but not a combination of these traits. Change in genotype frequencies due to selection The change in genotype frequencies due to selection follows immediately from the definition of relative fitness, . Thus, a genotype's frequency will decline or increase depending on whether its fitness is lower or greater than the mean fitness, respectively. In the particular case that there are only two genotypes of interest (e.g. representing the invasion of a new mutant allele), the change in genotype frequencies is often written in a different form. Suppose that two genotypes and have fitnesses and , and frequencies and , respectively. Then , and so . Thus, the change in genotype 's frequency depends crucially on the difference between its fitness and the fitness of genotype . Supposing that is more fit than , and defining the selection coefficient by , we obtain , where the last approximation holds for . In other words, the fitter genotype's frequency grows approximately logistically. History The British sociologist Herbert Spencer coined the phrase "survival of the fittest" in his 1864 work Principles of Biology to characterise what Charles Darwin had called natural selection. The British-Indian biologist J.B.S. Haldane was the first to quantify fitness, in terms of the modern evolutionary synthesis of Darwinism and Mendelian genetics starting with his 1924 paper A Mathematical Theory of Natural and Artificial Selection. The next further advance was the introduction of the concept of inclusive fitness by the British biologist W.D. Hamilton in 1964 in his paper on The Genetical Evolution of Social Behaviour. Genetic load Genetic load measures the average fitness of a population of individuals, relative either to a theoretical genotype of optimal fitness, or relative to the most fit genotype actually present in the population. Consider n genotypes , which have the fitnesses and the genotype frequencies respectively. Ignoring frequency-dependent selection, then genetic load () may be calculated as: Genetic load may increase when deleterious mutations, migration, inbreeding, or outcrossing lower mean fitness. Genetic load may also increase when beneficial mutations increase the maximum fitness against which other mutations are compared; this is known as the substitutional load or cost of selection.	Biology and health sciences	Basics_4	Biology
187897	https://en.wikipedia.org/wiki/Mastodon	Mastodon	A mastodon ( 'breast' + 'tooth') is a member of the genus Mammut (German for 'mammoth'), which was endemic to North America and lived from the late Miocene to the early Holocene. Mastodons belong to the order Proboscidea, the same order as elephants and mammoths (which belong to the family Elephantidae). Mammut is the type genus of the extinct family Mammutidae, which diverged from the ancestors of modern elephants at least 27–25 million years ago, during the Oligocene. Like other members of Mammutidae, the molar teeth of mastodons have zygodont morphology (where parallel pairs of cusps are merged into sharp ridges), which strongly differ from those of elephantids. In comparison to its likely ancestor Zygolophodon, Mammut is characterized by particularly long and upward curving upper tusks, reduced or absent tusks on the lower jaw, as well as the shortening of the mandibular symphysis (the frontmost part of the lower jaw), the latter two traits also having evolved in parallel separately in elephantids. Mastodons had an overall stockier skeletal build, a lower-domed skull, and a longer tail compared to elephantids. Fully grown male M. americanum are thought to have been have been to at shoulder height and from to in body mass on average. The size estimates suggest that American mastodon males were on average heavier than any living elephant species; they were typically larger than Asian elephants and African forest elephants of both sexes but shorter than male African bush elephants. M. americanum, known as an "American mastodon" or simply "mastodon," had a long and complex paleontological history spanning all the way back to 1705 when the first fossils were uncovered from Claverack, New York in the American colonies. Because of the uniquely shaped molars with no modern analogues in terms of large animals, the species caught wide attention of European researchers and influential Americans before and after the American Revolution to the point of, according to American historians Paul Semonin and Keith Stewart Thomson, bolstering American nationalism and contributing to a greater understanding of extinctions. Taxonomically, it was first recognized as a distinct species by Robert Kerr in 1792 then classified to its own genus Mammut by Johann Friedrich Blumenbach in 1799, thus making it amongst the first fossil mammal genera to be erected with undisputed taxonomic authority. The genus served as a wastebasket taxon for proboscidean species with superficially similar molar teeth morphologies but today includes 7 definite species, 1 of questionable affinities, and 4 other species from Eurasia that are pending reassessments to other genera. Mastodons are considered to have had a predominantly browsing-based diet on leaves, fruits, and woody parts of plants. This allowed mastodons to niche partition with other members of Proboscidea in North America, like gomphotheres and the Columbian mammoth, who had shifted to mixed feeding or grazing by the late Neogene-Quaternary. It is thought that mastodon behaviors were not much different from elephants and mammoths, with females and juveniles living in herds and adult males living largely solitary lives plus entering phases of aggression similar to the musth exhibited by modern elephants. Mammut achieved maximum species diversity in the Pliocene, though the genus is known from abundant fossil evidence in the Late Pleistocene. Mastodons for at least a few thousand years prior to their extinction coexisted with Paleoindians, who were the first humans to have inhabited North America. Evidence has been found that Paleoindians (including those of the Clovis culture) hunted mastodons based on the finding of mastodon remains with cut marks and/or with lithic artifacts. Mastodons disappeared along with many other North American animals, including most of its largest animals (megafauna), as part of the end-Pleistocene extinction event around the end of the Late Pleistocene-early Holocene, the causes typically being attributed to human hunting, severe climatic phases like the Younger Dryas, or some combination of the two. The American mastodon had its last recorded occurrence in the earliest Holocene around 11,000 years ago, which is considerably later than other North American megafauna species. Today, the American mastodon is one of the most well-known fossil species in both academic research and public perception, the result of its inclusion in American popular culture. Taxonomy Research history Earliest finds In a letter dating to 1713, Edward Hyde, 3rd Earl of Clarendon (known also as Lord Cornbury) from New York reported to the Royal Society learned society of Great Britain that in 1705, a large-sized tooth was found near the side of the Hudson River by a Dutch country-fellow and was sold to New York General Assembly member Van Bruggen for a gill of rum, and Bruggen eventually gave it to Cornbury. He then stated that he sent Johannis Abeel, a recorder of Albany, New York to dig near the original site of the tooth to find more bones. Abeel reported later that he went to the town of Claverack, New York where the original bones were found. American historian Paul Semonin said that the accounts written by Cornbury and Abeel match up with that written by in the July 30, 1705 entry in The Boston News-Letter. The account reported skeletal evidence of an antediluvian (or biblical) "giant" uncovered from Claverack. The femur and one of the teeth both dissolved before they could be further observed, however. Big Bone Lick In 1739, a French military expedition under the command of Charles III Le Moyne (known also as "Longueil") explored the locality of "Big Bone Lick" (located in what is now the US state of Kentucky) and gathered fossil bones and teeth there. The French naturalist Louis Jean-Marie Daubenton examined the fossil collection brought by Longueuil and compared it with specimens of extant elephants and Siberian mammoths in 1762. Daubenton said that the bones were discovered by Native Americans (probably Abenaki hunter-warriors). He came to the conclusion that the femur and tusk belonged to an elephant while the molars (or cheek teeth) came from a separate giant hippopotamus. In Shawnee tradition, the proboscideans roamed in herds and were hunted by giants, who both eventually died out. The accounts told by the Shawnee individuals in 1762 are the oldest known documented interpretations of the "Ohio" fossils, although the traditions may have had been told for generations. In 1767, Peter Collinson credited Irish trader George Croghan for having sent him and Benjamin Franklin fossil evidence of the mysterious proboscideans, using them for his studies. He concluded that the peculiar grinders (the molars) were built for herbivorous diets of branches of trees and shrubs as well as other vegetation, a view later followed by Franklin. In 1768, Scottish anatomist William Hunter recorded that he and his brother John Hunter observed that the teeth were not like those of modern elephants. He determined that the "grinders" from Ohio were of a carnivorous animal but believed that the tusks belonged to the same animal. After examining fossils from Franklin and Lord Shelburne, Hunter was convinced that the "pseudo-elephant", or "animal incognitum" (shortened as "incognitum"), was an animal species separate from elephants that might have also been the same as the proboscideans found in Siberia. He concluded his article with the opinion that although regrettable to philosophers, humanity should be thankful to heaven that the animal, if truly carnivorous, was extinct. Early American observations In 1785, Reverend Robert Annan wrote an account recalling an event in which workers discovered bones in his farm near the Hudson River in New York in fall of 1780. The workers found four molars in addition to another that was broken and thrown away. They also uncovered bones, including vertebrae that broke shortly thereafter. Annan expressed his confusion at what the animal could be but speculated based on its "grinders" that it was carnivorous in diet. He speculated also that it was probably extinct due to some catastrophe within the globe. American statesman Thomas Jefferson stated his thoughts on	Biology and health sciences	Proboscidea	Animals
187961	https://en.wikipedia.org/wiki/Acela	Acela	The Acela ( ; originally the Acela Express until September 2019) is Amtrak's flagship passenger train service along the Northeast Corridor (NEC) in the Northeastern United States between Washington, D.C. and Boston via 13 intermediate stops, including Baltimore, New York City and Philadelphia. Acela trains are the fastest in the Americas, reaching (qualifying as high-speed rail), but only over of the route. Acela carried more than 2.9 million passengers in fiscal year 2023, second only to the slower and less expensive Northeast Regional, which had over 9 million passengers. Ridership is down from the pre-Covid-19 pandemic high of 3,557,455 passengers in 2019. Its 2016 revenue of $585 million was 25% of Amtrak's total. Acela operates along routes that are used by freight and slower regional passenger traffic, and reaches the maximum allowed speed of the tracks only along some sections, with the fastest peak speed along segments between Mansfield, Massachusetts, and Richmond, Rhode Island, and South Brunswick and Trenton, New Jersey. Acela trains use active tilting technology, which helps control lateral centrifugal force, allowing the train to travel at higher speeds on the sharply curved NEC without disturbing passengers. The high-speed operation occurs mostly along the route from Pennsylvania Station in New York City to Union Station in Washington, D.C., with a fastest scheduled time of 2 hours and 45 minutes and an average speed of , including time spent at intermediate stops. Over this route, Acela and the Northeast Regional service captured an 83% share of air/train commuters between New York and Washington in 2021, up from 37% in 2000. The Acela speed is limited by traffic and infrastructure on the route's northern half. On the section from Boston's South Station to New York's Penn Station, the fastest scheduled time is 3 hours and 30 minutes, or an average speed of . Along this section, Acela has captured a 54% share of the combined train and air market. The entire route from Boston to Washington takes between 6 hours, 38 minutes and 6 hours, 50 minutes, at an average speed of around . The present Acela Express equipment will be replaced by new Avelia Liberty trainsets in early 2025. The new trains will have greater passenger capacity and an enhanced active tilt system that will allow higher speed on the many curved sections of the route. History Background Following the success of Japan's newly inaugurated Shinkansen network, the High Speed Ground Transportation Act of 1965 authorized the U.S. government to explore the creation of high-speed rail, which resulted in the introduction of the higher-speed Metroliner trains between Washington, D.C., and New York City in 1969, the predecessor to Acela. During the 1980s, the U.S. Federal Railroad Administration explored the possibilities of high-speed rail in the United States. On December 18, 1991, five potential high speed rail corridors were authorized, including the Northeast Corridor. In the early 1990s, Amtrak tested several different high-speed trains from Europe on the Northeast Corridor. An X 2000 train was leased from Sweden for test runs from October 1992 to January 1993, followed by revenue service between Washington, D.C., and New York City from February to May and August to September 1993. Siemens showed the ICE 1 train from Germany, organizing the ICE Train North America Tour which started to operate on the Northeast Corridor on July 3, 1993. Building and development With the testing of the trains from Europe complete, Amtrak was able to define a set of specifications for high-speed equipment and in October 1994, Amtrak requested bids from train manufacturers for a trainset that could reach . A consortium of Bombardier (75%) and GEC Alsthom (now Alstom) (25%) was selected in March 1996. On March 9, 1999, Amtrak unveiled its plan for the Acela Express, a high-speed train on the Northeast Corridor between Washington, D.C., and Boston. Several changes were made to the corridor to make it suitable for higher-speed electric trains. The Northend Electrification Project extended existing electrification from New Haven, Connecticut, to Boston to complete the overhead power supply along the route, and several grade crossings were improved or removed. Prior to 2000, all trains bound for Boston had to switch to diesel power at New Haven. A pilot trainset was completed by early 2000 and sent to Transportation Technology Center (TTC) for testing in June 2000. An inaugural VIP run of the Acela occurred on November 16, 2000, with the VIP train being led by power car number 2020 with no. 2009 at the opposite end, followed by the first revenue run on December 11, 2000, a few months after the intended date. Cost Amtrak's original contract with the Bombardier-Alstom consortium was for the delivery of 20 trainsets (six coaches each, with power cars at front and rear) for $800 million. By 2004, Amtrak had settled contract disputes with the consortium, paying a total of $1.2 billion for the 20 trainsets plus 15 extra high-speed locomotives and the construction of maintenance facilities in Boston, New York, and Washington. Impact of the Acela By 2005, Amtrak's share of the common-carrier market between New York and Boston had reached 40%, from 18% pre-Acela. With the increasing popularity of the faster, modern Acela Express, Metroliner service was phased out in late 2006. To meet the demand, more Acela services were added in September 2005. By August 2008 crowding had become noticeable. By 2011, the Acela fleet had reached half of its designed service life. Amtrak proposed several replacement options, including one as part of its A Vision for High-Speed Rail in the Northeast Corridor. In 2011, Amtrak announced that forty new Acela coaches would be ordered in 2012 to increase capacity on existing trainsets. The existing trains would have received two more coaches, lengthening the trainsets from a 1-6-1 configuration to 1-8-1 (power car—passenger cars—power car). The longer trainsets would have required the modifications of the Acela maintenance facilities in Boston, New York and Washington. The first of the stretched trainsets was to have entered service in fiscal year 2014. This plan was cancelled in 2012 in favor of replacing, rather than refurbishing, the Acela fleet. In January 2014, Amtrak issued a request for proposals on 28 or more new model Acela trainsets, in a combined order with the California High-Speed Rail Authority. These bids were due May 17, 2014. After discussions with manufacturers, Amtrak and the California High Speed Rail Authority concluded their needs were too disparate for common rolling stock and decided not to pursue the joint option. Branding Before the introduction of the Acela, there were several classes of trains on the Northeast Corridor: the express Metroliners, the Philadelphia-New York Clockers, Empire Service trains between New York City and Niagara Falls via the Empire Corridor, Keystone Service trains between New York City and to Harrisburg via the Keystone Corridor, and the umbrella term NortheastDirect, applied to other trains on the corridor (in addition to unique names assigned to many departures). The Acela name was announced on March 9, 1999, as a part of the original announcement of the service itself. The branding team based the name "Acela" on the ideas of acceleration and excellence. At the same time, Amtrak launched what it called the Capstone Program, a short-lived plan to rebrand the NortheastDirect, Keystone Service and Empire Service trains as Acela Regional and the Clocker trains as Acela Commuter. The Acela Regional name was first applied to NortheastDirect trains 130–133 on January 31, 2000. Those trains were the first electrified trains to run on the full Northeast Corridor between Washington, D.C., and Boston. As more trains were electrified, they too were rebranded. Following mass rider confusion between the three services, the branding was removed from the lower-speed Acela Regional and Acela Commuter trains in 2003. On September 23, 2019, Amtrak shortened the name of the service from Acela Express to simply Acela. At the same time, Amtrak introduced the Acela Nonstop, a direct train from Washington, D.C., to New York's Penn Station. The nonstop service was temporarily suspended as of March 10, 2020, due to low ridership caused by the COVID-19 pandemic. Engineering First-generation trainsets The first-generation Acela trainset is a unique set of vehicles designed specifically to satisfy governmental rolling stock requirements established primarily by the Federal Railroad Administration (FRA). This includes the ability to withstand a collision with a freight train at speed without collapsing. Most manufacturers that bid on the Acela were unable to meet the structural requirements, due to increased costs and complications for the manufacture of the trains, and the need for manufacturers to make significant engineering changes to their standard designs. In the end, only three qualified bidders remained: ABB (Swedish-Swiss manufacturer of the X 2000 train), Siemens (manufacturer of the German ICE), and a consortium of Bombardier (manufacturer of the LRC trains) and Alstom (manufacturer of the French TGV). The design, using identical power cars at each end which operate on voltages of 12 kV, 12.5 kV, and 25 kV AC, and either 25 or 60 Hz frequency, derives several components from the TGV, such as the third-generation TGV's traction system (including the four asynchronous AC motors per power car, rectifiers, inverters, and regenerative braking), the trucks/bogies structure (a long-wheelbase dual transom H frame welded steel with outboard mounted tapered roller bearings), the brake discs (although there are only three per axle, versus four on the TGV), and crash energy management techniques to control structural deformation in the event of an accident. The tilting carriages are based upon Bombardier's earlier LRC trains used on Via Rail rather than the TGV's non-tilting articulated trailers. Acela power cars and passenger cars are much heavier than those of the TGV in order to meet the FRA's crash standards. French and Canadian crews testing the Acela referred to it as "the pig" due to its weight. The extra weight leads to the Acela power-to-weight ratio being about 22.4 hp per tonne, compared to 30.8 hp for a SNCF TGV Reseau trainset. The Tier II crash standards, adopted in 1999, have also resulted in the passenger cars being designed without steps and trapdoors, which means that the trainsets can only serve lines with high-level platforms such as the Northeast Corridor. Acela trains are semi-permanently coupled (but not articulated as in the TGV) and are referred to as trainsets. Bombardier later used the Acela carriage design and a diesel/gas turbine variant of the power car for its experimental JetTrain. Second-generation trainsets On August 26, 2016, then-Vice President Joe Biden announced a $2.45 billion federal loan package to pay for new equipment for the Acela Express service, as well as upgrades to the NEC. The loans will finance 28 Avelia Liberty trainsets that will be built by Alstom in Hornell and Rochester, New York, and will replace the existing fleet of twenty Acela trainsets. The fleet expansion will allow for hourly New York-Boston service all day and half-hourly New York-Washington service at peak hours. The new trainsets will be longer, have 386 seats compared to 304 on Acela Express (a 27% increase) and will feature active tilt technology that will initially allow service to operate at 160 mph (260 km/h) and would allow for service if proposed infrastructure improvements are completed. The new trains were expected to be phased in between 2021 and 2022, after which the current fleet was to be retired. Trains are now expected to enter passenger service in spring 2025. Operating speeds Although the first-generation Acela Express trainsets were designed with a top speed of and the second-generation Avelia Liberty trainsets will be designed to reach , the existing infrastructure of the Northeast Corridor significantly limits speeds. The maximum speed limit on the Northeast Corridor is on of the route, in four sections of track in Rhode Island, Massachusetts, and New Jersey. The Acela achieves an average speed (including stops) of between Washington and New York, and an average speed of from New York to Boston. The average speed over the entire route is . Speeds are limited by the route the corridor takes through urban areas, and there are several speed restrictions below over bridges or through tunnels that are over a century old. Altogether, Amtrak has identified 224 bridges along Acela route that are beyond their design life. South of the Delaware River, the Acela's top speed is . One limiting factor is the overhead catenary support system which was constructed before 1935 and lacks the constant-tension features of the new catenary east of New Haven. The Pennsylvania Railroad ran Metroliner test trains in the late 1960s as fast as and briefly intended to run the Metroliner service at speeds reaching . Certification testing for commercial operation at involving test runs at up to began between Trenton and New Brunswick in September 2012. Passenger operation at began in this region in late May 2022. The fastest schedule between New York and Washington, DC was 2 hours, 43 minutes in 2012. $450 million was allotted by President Barack Obama's administration to replace catenary and upgrade signals between Trenton and New Brunswick, which will allow speeds of over a stretch. The improvements were scheduled to be completed in 2016, but have been delayed; the project was partially completed in late May 2022, with the remainder projected in 2024. This section of track holds the record for the highest speed by a train in the US, which is , achieved in a test run by the U.S./Canada-built UAC TurboTrain on December 20, 1967. North of New York City, Amtrak upgraded the track along the Connecticut shoreline east of New Haven to allow maximum speeds in excess of , in preparation for the Acela launch. Although this area contains the fastest current operating speeds (), it also has the slowest section of the NEC: between New Rochelle, New York, and New Haven, Connecticut. This section is owned by Metro-North Railroad and the Connecticut Department of Transportation and is heavily used by commuter trains which limit the speed of the Acela. Amtrak's trains achieve only on a limited stretch in New York State and rarely exceed at any time eastbound through Connecticut until reaching New Haven. In 1992, ConnDOT began plans to upgrade the catenary system and replace outdated bridges on the New Haven Line to enable the Acela to run slightly faster. the catenary replacement and bridge work were under way and expected to be completed by mid-2018. On July 9, 2007, Amtrak introduced a limited-stop round trip, with trains stopping only at Philadelphia between New York and Washington. This shortened the trip between the two cities to 2 hours 35 minutes, making the trip roughly an hour faster than some of the Northeast Regional train services. These trains were an experiment to find ways to expedite travel time on the Acela; Amtrak has since dropped them. High speed infrastructure The dense population of the northeastern United States makes the Northeast Corridor the most heavily traveled portion of the American passenger rail system. Two-thirds of rail passengers in the United States live in or near New York City, also home to the nation's busiest passenger rail station, Penn Station. In order to compete with airliners, Amtrak needed to increase the speed of trains in the region. The former Shore Line from New Haven to Boston is burdened by sharp turns and grade crossings, the crossings being of special concern. Tilting enables passengers to ride more comfortably on curved sections of track faster than would otherwise be possible, by leaning into the bend. Acela trainsets use active tilting above on most of the system, but some segments of track in the Northeast Corridor are too close together for the cars to safely tilt while maintaining FRA minimum space between trains on parallel tracks. Metro-North Railroad restricts tilting on the segment of track north of New York which it owns. The system was originally designed for a 6.8° tilt, but the cars were redesigned wider to accommodate wider seats and aisles that reduced allowable tilt to 4.2° to fit within the clearance constraints of the existing tracks. Traveling at higher than also requires constant-tension catenary, which is only implemented on the more modern catenary system north of New York City. South of New York City, the trains are restricted to . By comparison, the Northeast Regional and the now-defunct Metroliner service reached . Acela service was originally expected to begin in late 1999 but was delayed. The catenary system could not support the intended speeds between Washington DC and New York City, but the newer system between New York City and Boston allows the higher speeds. Attention was drawn to the decreased 4.2° tilt, but this was not the root of the speed problem, as the tracks from New York to Boston are similar to those between New York and Washington, and the tilt mechanism is not the factor enabling higher speeds. Following repairs, the first Acela service began on December 11, 2000, a year behind schedule. Acela travels between Boston and New York in about three and a half hours (an improvement of half an hour); New York to Washington runs take a minimum of two hours and forty-five minutes. These schedules, as well as the relative convenience of direct downtown-to-downtown rail service as opposed to air travel, especially after the September 11 attacks, have made the Acela Express more competitive with the air shuttles. Due to this competition, Southwest Airlines canceled service between Washington and New York. Platform track speeds Due to the high speed at which Acela trains bypass platforms of local stations, concerns have mounted in some communities over inadequate warnings and safeguards for passengers waiting for other trains, including that the two-foot wide yellow platform markings may not keep people at a safe distance. At Kingston station in Rhode Island and Mansfield station in Massachusetts, Acela trains pass by at . Suggestions include platform safety barriers, or use of different announcements for approaching Acela trains versus slower ones. Outages In August 2002, shortly after their introduction, Acela trainsets were briefly removed from service when the brackets that connected truck (bogie) dampers (shocks) to the powerunit carbodies ("yaw dampers") were found to be cracking. The Acela returned to service when a program of frequent inspections was instituted. The damper brackets have since been redesigned and old brackets replaced by the newer design. On April 15, 2005, the Acela was removed from service when cracks were found in the disc brakes of many passenger coaches. The Bombardier-Alstom consortium replaced the discs under warranty. Limited service resumed in July 2005, as a portion of the fleet operated with new brake discs. Metroliner trains, which the Acela Express was intended to replace, filled in during the outage. Amtrak announced on September 21, 2005, that all 20 trainsets had been returned to full operation. In October 2012, Acela service was cancelled immediately before, during, and after Hurricane Sandy, which damaged the North River Tunnels causing lasting delays and reliability problems. In March 2020, all Acela trips were suspended as part of a round of service reduction in response to the COVID-19 pandemic in the United States. Amtrak resumed Acela service on June 1, 2020. Service Composition The production sets are formed as follows: The Acela Express trainset consists of two power cars, a Café car, a First Class car, and four Business Class cars, semi-permanently coupled together. It has fewer seats than its regional service counterparts. The First Class car has 44 seats, being three seats across (one on one side, two on the other side), four-seat tables and assigned seating. There are 260 Business Class seats on each trainset; these cars have four seats across (two on each side), four-seat tables, and assigned seating. Baggage may be stowed in overhead compartments or underneath seats. Trains are wheelchair-accessible. Each car has one or two toilets, with one being ADA compliant. The Business Class car adjacent to First Class is designated as a quiet car, where passengers are asked to refrain from loud talking and phone conversations. Automatic sliding doors between cars reduce noise. Operations and staffing Acela offers two classes of seating, Business Class and First Class. Unlike most other Amtrak trains, Business Class is the de facto standard class on Acela trains; there is no coach service. Acela maintenance is generally taken care of at the Ivy City facility in Washington, DC; Sunnyside Yard in Queens, New York; or Southampton Street Yard in Boston. The Acela trainsets underwent minor refurbishments between mid-2009 and 2010 at Penn Coach Yard, next to 30th Street Station in Philadelphia, Pennsylvania. These refurbishments included new blue leather seats throughout the trainset. In May 2018, Amtrak announced a 14-month program to refresh the interiors of the Acela trainsets, including new seat cushions and covers, new aisle carpeting, and a deep clean. This refurbishment program has been completed as of June 2019. Wi-Fi service Wireless Internet station service began in 2004. In 2010, with services provided by The GBS Group, all Acela trains began offering "AmtrakConnect" supporting IEEE 802.11a/b/g/n, 2.4 GHz and 5 GHz and standard VPN connections. In 2016, Amtrak upgraded to a faster wifi service. Staffing Generally, Amtrak train crews consist of an engineer, a conductor, and at least one assistant conductor. Acela trains also have an On-Board Service crew consisting of two First Class attendants and a Café Car attendant. In addition to the food service provided in the Café Car, on most trains an attendant will also provide at-seat cart service, serving refreshments throughout the train. First Class passengers are served meals at their seats on all services. Notable incidents During the Northeast blackout of 2003, a northbound Acela Express train was stuck on the Hell Gate Bridge for over nine hours, until a rescue engine from Sunnyside Yard was able to tow the train back to New York's Penn Station. The first Acela grade crossing accident occurred on September 27, 2005, when a car rolled under closed crossing gate arms in Waterford, Connecticut, and was struck by a train traveling at , killing three automobile passengers. None of the 130 Acela passengers were injured. The gates were found to have been functioning properly, but the incident drew much criticism regarding the eleven remaining grade crossings along Amtrak's busy Northeast Corridor. On March 24, 2017, an Acela Express train derailed at low speed in New York's Penn Station, during morning rush hour. All 248 passengers were safely evacuated. The derailment was caused by a defective section of track, of which Amtrak was aware, but had not yet fixed. On February 6, 2018, Acela Express train No. 2150 split apart between the first and second cars in the trainset, at , near Havre de Grace, Maryland. There were no injuries of the crew nor the 52 passengers on board, who were transferred to Northeast Regional train No. 180. Station stops A limited number of Acela trains previously stopped at New Rochelle, New York; New London, Connecticut; and Trenton, New Jersey; service was eliminated in 2021, 2022 and 2023, respectively.	Technology	High-speed rail	null
188037	https://en.wikipedia.org/wiki/Pnictogen	Pnictogen	\|- ! colspan=2 style="text-align:left;" \| ↓ Period \|- ! 2 \| \|- ! 3 \| \|- ! 4 \| \|- ! 5 \| \|- ! 6 \| \|- ! 7 \| \|- \| colspan="2"\| Legend \|} The pnictogens ( or ; from "to choke" and -gen, "generator") are the chemical elements in group 15 of the periodic table. This group is also known as the nitrogen group or nitrogen family. Group 15 consists of the elements nitrogen (N), phosphorus (P), arsenic (As), antimony (Sb), bismuth (Bi), and moscovium (Mc). Since 1988, it has been called Group 15 by the IUPAC. Before that, in America it was called Group VA, owing to a text by H. C. Deming and the Sargent-Welch Scientific Company, while in Europe it was called Group VB, which the IUPAC had recommended in 1970. (Pronounced "group five A" and "group five B"; "V" is the Roman numeral 5). In semiconductor physics, it is still usually called Group V. The "five" ("V") in the historical names comes from the "pentavalency" of nitrogen, reflected by the stoichiometry of compounds such as N2O5. They have also been called the pentels. Characteristics Chemical Like other groups, the members of this family manifest similar patterns in electron configuration, notably in their valence shells, resulting in trends in chemical behavior. This group has a defining characteristic whereby each component element has 5 electrons in its valence shell, that is, 2 electrons in the s sub-shell and 3 unpaired electrons in the p sub-shell. They are therefore 3 electrons shy of filling their valence shell in their non-ionized state. The Russell-Saunders term symbol of the ground state in all elements in the group is 4S. The most important elements of this group to life on Earth are nitrogen (N), which in its diatomic form is the principal component of air, and phosphorus (P), which, like nitrogen, is essential to all known forms of life. Compounds Binary compounds of the group can be referred to collectively as pnictides. Magnetic properties of pnictide compounds span the cases of diamagnetic systems (such as BN or GaN) and magnetically ordered systems (MnSb is paramagnetic at elevated temperatures and ferromagnetic at room temperature); the former compounds are usually transparent and the latter metallic. Other pnictides include the ternary rare-earth (RE) main-group variety of pnictides. These are in the form of , where M is a carbon group or boron group element and Pn is any pnictogen except nitrogen. These compounds are between ionic and covalent compounds and thus have unusual bonding properties. These elements are also noted for their stability in compounds due to their tendency to form covalent double bonds and triple bonds. This property of these elements leads to their potential toxicity, most evident in phosphorus, arsenic, and antimony. When these substances react with various chemicals of the body, they create strong free radicals that are not easily processed by the liver, where they accumulate. Paradoxically, this same strong bonding causes nitrogen's and bismuth's reduced toxicity (when in molecules), because these strong bonds with other atoms are difficult to split, creating very unreactive molecules. For example, , the diatomic form of nitrogen, is used as an inert gas in situations where using argon or another noble gas would be too expensive. Formation of multiple bonds is facilitated by their five valence electrons, as the octet rule permits a pnictogen to accept three electrons on covalent bonding. As 5 3, it leaves two unused electrons in a lone pair unless there is a positive charge around (like in ). When a pnictogen forms only three single bonds, effects of the lone pair typically results in trigonal pyramidal molecular geometry. Oxidation states The light pnictogens (nitrogen, phosphorus, and arsenic) tend to form −3 charges when reduced, completing their octet. When oxidized or ionized, pnictogens typically take an oxidation state of +3 (by losing all three p-shell electrons in the valence shell) or +5 (by losing all three p-shell and both s-shell electrons in the valence shell). However heavier pnictogens are more likely to form the +3 oxidation state than lighter ones due to the s-shell electrons becoming more stabilized. −3 oxidation state Pnictogens can react with hydrogen to form pnictogen hydrides such as ammonia. Going down the group, to phosphane (phosphine), arsane (arsine), stibane (stibine), and finally bismuthane (bismuthine), each pnictogen hydride becomes progressively less stable (more unstable), more toxic, and has a smaller hydrogen-hydrogen angle (from 107.8° in ammonia to 90.48° in bismuthane). (Also, technically, only ammonia and phosphane have the pnictogen in the −3 oxidation state because, for the rest, the pnictogen is less electronegative than hydrogen.) Crystal solids featuring pnictogens fully reduced include yttrium nitride, calcium phosphide, sodium arsenide, indium antimonide, and even double salts like aluminum gallium indium phosphide. These include III-V semiconductors, including gallium arsenide, the second-most widely used semiconductor after silicon. +3 oxidation state Nitrogen forms a limited number of stable III compounds. Nitrogen(III) oxide can only be isolated at low temperatures, and nitrous acid is unstable. Nitrogen trifluoride is the only stable nitrogen trihalide, with nitrogen trichloride, nitrogen tribromide, and nitrogen triiodide being explosive—nitrogen triiodide being so shock-sensitive that the touch of a feather detonates it (the last three actually feature nitrogen in the −3 oxidation state). Phosphorus forms a +III oxide which is stable at room temperature, phosphorous acid, and several trihalides, although the triiodide is unstable. Arsenic forms +III compounds with oxygen as arsenites, arsenous acid, and arsenic(III) oxide, and it forms all four trihalides. Antimony forms antimony(III) oxide and antimonite but not oxyacids. Its trihalides, antimony trifluoride, antimony trichloride, antimony tribromide, and antimony triiodide, like all pnictogen trihalides, each have trigonal pyramidal molecular geometry. The +3 oxidation state is bismuth's most common oxidation state because its ability to form the +5 oxidation state is hindered by relativistic properties on heavier elements, effects that are even more pronounced concerning moscovium. Bismuth(III) forms an oxide, an oxychloride, an oxynitrate, and a sulfide. Moscovium(III) is predicted to behave similarly to bismuth(III). Moscovium is predicted to form all four trihalides, of which all but the trifluoride are predicted to be soluble in water. It is also predicted to form an oxychloride and oxybromide in the +III oxidation state. +5 oxidation state For nitrogen, the +5 state is typically serves as only a formal explanation of molecules like N2O5, as the high electronegativity of nitrogen causes the electrons to be shared almost evenly. Pnictogen compounds with coordination number 5 are hypervalent. Nitrogen(V) fluoride is only theoretical and has not been synthesized. The "true" +5 state is more common for the essentially non-relativistic typical pnictogens phosphorus, arsenic, and antimony, as shown in their oxides, phosphorus(V) oxide, arsenic(V) oxide, and antimony(V) oxide, and their fluorides, phosphorus(V) fluoride, arsenic(V) fluoride, antimony(V) fluoride. They also form related fluoride-anions, hexafluorophosphate, hexafluoroarsenate, hexafluoroantimonate, that function as non-coordinating anions. Phosphorus even forms mixed oxide-halides, known as oxyhalides, like phosphorus oxychloride, and mixed pentahalides, like phosphorus trifluorodichloride. Pentamethylpnictogen(V) compounds exist for arsenic, antimony, and bismuth. However, for bismuth, the +5 oxidation state becomes rare due to the relativistic stabilization of the 6s orbitals known as the inert-pair effect, so that the 6s electrons are reluctant to bond chemically. This causes bismuth(V) oxide to be unstable and bismuth(V) fluoride to be more reactive than the other pnictogen pentafluorides, making it an extremely powerful fluorinating agent. This effect is even more pronounced for moscovium, prohibiting it from attaining a +5 oxidation state. Other oxidation states Nitrogen forms a variety of compounds with oxygen in which the nitrogen can take on a variety of oxidation states, including +II, +IV, and even some mixed-valence compounds and very unstable +VI oxidation state. In hydrazine, diphosphane, and organic derivatives of the two, the nitrogen or phosphorus atoms have the −2 oxidation state. Likewise, diimide, which has two nitrogen atoms double-bonded to each other, and its organic derivatives have nitrogen in the oxidation state of −1. Similarly, realgar has arsenic–arsenic bonds, so the arsenic's oxidation state is +II. A corresponding compound for antimony is Sb2(C6H5)4, where the antimony's oxidation state is +II. Phosphorus has the +1 oxidation state in hypophosphorous acid and the +4 oxidation state in hypophosphoric acid. Antimony tetroxide is a mixed-valence compound, where half of the antimony atoms are in the +3 oxidation state, and the other half are in the +5 oxidation state. It is expected that moscovium will have an inert-pair effect for both the 7s and the 7p1/2 electrons, as the binding energy of the lone 7p3/2 electron is noticeably lower than that of the 7p1/2 electrons. This is predicted to cause +I to be a common oxidation state for moscovium, although it also occurs to a lesser extent for bismuth and nitrogen. Physical The pnictogens exemplify the transition from nonmetal to metal going down the periodic table: a gaseous diatomic nonmetal (N), two elements displaying many allotropes of varying conductivities and structures (P and As), and then at least two elements that only form metallic structures in bulk (Sb and Bi; probably Mc as well). All the elements in the group are solids at room temperature, except for nitrogen which is gaseous at room temperature. Nitrogen and bismuth, despite both being pnictogens, are very different in their physical properties. For instance, at STP nitrogen is a transparent non-metallic gas, while bismuth is a silvery-white metal. The densities of the pnictogens increase towards the heavier pnictogens. Nitrogen's density is 0.001251 g/cm3 at STP. Phosphorus's density is 1.82 g/cm3 at STP, arsenic's is 5.72 g/cm3, antimony's is 6.68 g/cm3, and bismuth's is 9.79 g/cm3. Nitrogen's melting point is −210 °C and its boiling point is −196 °C. Phosphorus has a melting point of 44 °C and a boiling point of 280 °C. Arsenic is one of only two elements to sublimate at standard pressure; it does this at 603 °C. Antimony's melting point is 631 °C and its boiling point is 1587 °C. Bismuth's melting point is 271 °C and its boiling point is 1564 °C. Nitrogen's crystal structure is hexagonal. Phosphorus's crystal structure is cubic. Arsenic, antimony, and bismuth all have rhombohedral crystal structures. Nuclear All pnictogens up to antimony have at least one stable isotope; bismuth has no stable isotopes, but has a primordial radioisotope with a half-life much longer than the age of the universe (209Bi); and all known isotopes of moscovium are synthetic and highly radioactive. In addition to these isotopes, traces of 13N, 32P, and 33P occur in nature, along with various bismuth isotopes (other than 209Bi) in the decay chains of thorium and uranium. History The nitrogen compound sal ammoniac (ammonium chloride) has been known since the time of the Ancient Egyptians. In the 1760s two scientists, Henry Cavendish and Joseph Priestley, isolated nitrogen from air, but neither realized the presence of an undiscovered element. It was not until several years later, in 1772, that Daniel Rutherford realized that the gas was indeed nitrogen. The alchemist Hennig Brandt first discovered phosphorus in Hamburg in 1669. Brandt produced the element by heating evaporated urine and condensing the resulting phosphorus vapor in water. Brandt initially thought that he had discovered the Philosopher's Stone, but eventually realized that this was not the case. Arsenic compounds have been known for at least 5000 years, and the ancient Greek Theophrastus recognized the arsenic minerals called realgar and orpiment. Elemental arsenic was discovered in the 13th century by Albertus Magnus. Antimony was well known to the ancients. A 5000-year-old vase made of nearly pure antimony exists in the Louvre. Antimony compounds were used in dyes in the Babylonian times. The antimony mineral stibnite may have been a component of Greek fire. Bismuth was first discovered by an alchemist in 1400. Within 80 years of bismuth's discovery, it had applications in printing and decorated caskets. The Incas were also using bismuth in knives by 1500. Bismuth was originally thought to be the same as lead, but in 1753, Claude François Geoffroy proved that bismuth was different from lead. Moscovium was successfully produced in 2003 by bombarding americium-243 atoms with calcium-48 atoms. Names and etymology The term "pnictogen" (or "pnigogen") is derived from the ancient Greek word () meaning "to choke", referring to the choking or stifling property of nitrogen gas. It can also be used as a mnemonic for the two most common members, P and N. The term "pnictogen" was suggested by the Dutch chemist Anton Eduard van Arkel in the early 1950s. It is also spelled "pnicogen" or "pnigogen". The term "pnicogen" is rarer than the term "pnictogen", and the ratio of academic research papers using "pnictogen" to those using "pnicogen" is 2.5 to 1. It comes from the Greek root (choke, strangle), and thus the word "pnictogen" is also a reference to the Dutch and German names for nitrogen ( and , respectively, "suffocating substance": i.e., substance in air, unsupportive of breathing). Hence, "pnictogen" could be translated as "suffocation maker". The word "pnictide" also comes from the same root. Previously, the name pentels (from Greek , , five) was also used for this group. Occurrence Nitrogen makes up 25 parts per million of the Earth's crust, 5 parts per million of soil on average, 100 to 500 parts per trillion of seawater, and 78% of dry air. Most nitrogen on Earth is in nitrogen gas, but some nitrate minerals exist. Nitrogen makes up 2.5% of a typical human by weight. Phosphorus is 0.1% of the earth's crust, making it the 11th most abundant element. Phosphorus comprises 0.65 parts per million of soil and 15 to 60 parts per billion of seawater. There are 200 Mt of accessible phosphates on earth. Phosphorus makes up 1.1% of a typical human by weight. Phosphorus occurs in minerals of the apatite family, which are the main components of the phosphate rocks. Arsenic constitutes 1.5 parts per million of the Earth's crust, making it the 53rd most abundant element. The soils hold 1 to 10 parts per million of arsenic, and seawater carries 1.6 parts per billion of arsenic. Arsenic comprises 100 parts per billion of a typical human by weight. Some arsenic exists in elemental form, but most arsenic is found in the arsenic minerals orpiment, realgar, arsenopyrite, and enargite. Antimony makes up 0.2 parts per million of the earth's crust, making it the 63rd most abundant element. The soils contain 1 part per million of antimony on average, and seawater contains 300 parts per trillion on average. A typical human has 28 parts per billion of antimony by weight. Some elemental antimony occurs in silver deposits. Bismuth makes up 48 parts per billion of the earth's crust, making it the 70th most abundant element. The soils contain approximately 0.25 parts per million of bismuth, and seawater contains 400 parts per trillion of bismuth. Bismuth most commonly occurs as the mineral bismuthinite, but bismuth also occurs in elemental form or sulfide ores. Moscovium is a synthetic element which does not occur naturally. Production Nitrogen Nitrogen can be produced by fractional distillation of air. Phosphorus The principal method for producing phosphorus is to reduce phosphates with carbon in an electric arc furnace. Arsenic Most arsenic is prepared by heating the mineral arsenopyrite in the presence of air. This forms As4O6, from which arsenic can be extracted via carbon reduction. However, it is also possible to make metallic arsenic by heating arsenopyrite at 650 to 700 °C without oxygen. Antimony With sulfide ores, the method by which antimony is produced depends on the amount of antimony in the raw ore. If the ore contains 25% to 45% antimony by weight, then crude antimony is produced by smelting the ore in a blast furnace. If the ore contains 45% to 60% antimony by weight, antimony is obtained by heating the ore, also known as liquidation. Ores with more than 60% antimony by weight are chemically displaced with iron shavings from the molten ore, resulting in impure metal. If an oxide ore of antimony contains less than 30% antimony by weight, the ore is reduced in a blast furnace. If the ore contains closer to 50% antimony by weight, the ore is instead reduced in a reverberatory furnace. Antimony ores with mixed sulfides and oxides are smelted in a blast furnace. Bismuth Bismuth minerals do occur, in particular in the form of sulfides and oxides, but it is more economic to produce bismuth as a by-product of the smelting of lead ores or, as in China, of tungsten and zinc ores. Moscovium Moscovium is produced a few atoms at a time in particle accelerators by firing a beam of calcium-48 ions at americium-243 until the nuclei fuse. Applications Liquid nitrogen is a commonly used cryogenic liquid. Nitrogen in the form of ammonia is a nutrient critical to most plants' survival. Synthesis of ammonia accounts for about 1–2% of the world's energy consumption and the majority of reduced nitrogen in food. Phosphorus is used in matches and incendiary bombs. Phosphate fertilizer helps feed much of the world. Arsenic was historically used as a Paris green pigment, which has since been discontinued due to its extreme toxicity. Arsenic in the form of organoarsenic compounds is sometimes used in chicken feed. Antimony is alloyed with lead to produce some bullets. Antimony currency was briefly used in the 1930s in parts of China, but was discontinued as antimony is both soft and toxic. Bismuth subsalicylate is the active ingredient in Pepto-Bismol. Bismuth chalcogenides are being studied in cancerous mice as a candidate for use in improving radiation therapy in human cancer patients. Moscovium is too unstable and scarce to have any known practical application. Biological role Nitrogen is a component of molecules critical to life on earth, such as DNA and amino acids. Nitrates occur in some plants, due to bacteria present in the nodes of the plant. This is seen in leguminous plants such as peas or spinach and lettuce. A typical 70 kg human contains 1.8 kg of nitrogen. Phosphorus in the form of phosphates occur in compounds important to life, such as DNA and ATP. Humans consume approximately 1 g of phosphorus per day. Phosphorus is found in foods such as fish, liver, turkey, chicken, and eggs. Phosphate deficiency is a problem known as hypophosphatemia. A typical 70 kg human contains 480 g of phosphorus. Arsenic promotes growth in chickens and rats, and may be essential for humans in small quantities. Arsenic has been shown to be helpful in metabolizing the amino acid arginine. There are 7 mg of arsenic in a typical 70 kg human. Antimony is not known to have a biological role. Plants take up only trace amounts of antimony. There are approximately 2 mg of antimony in a typical 70 kg human. Bismuth is not known to have a biological role. Humans ingest on average less than 20 μg of bismuth per day. There is less than 500 μg of bismuth in a typical 70 kg human. Moscovium is too unstable to occur in nature or have a known biological role. Moscovium does not typically occur in organisms in any meaningful amount. Toxicity Nitrogen gas is completely non-toxic, but breathing in pure nitrogen gas is deadly, because it causes nitrogen asphyxiation. The build-up of nitrogen bubbles in the blood, such as those that may occur during scuba diving, can cause a condition known as the "bends" (decompression sickness). Many nitrogen compounds such as hydrogen cyanide and nitrogen-based explosives are also highly dangerous. White phosphorus, an allotrope of phosphorus, is toxic, with 1 mg per kg bodyweight being a lethal dose. White phosphorus usually kills humans within a week of ingestion by attacking the liver. Breathing in phosphorus in its gaseous form can cause an industrial disease called "phossy jaw", which eats away the jawbone. White phosphorus is also highly flammable. Some organophosphorus compounds can fatally block certain enzymes in the human body. Elemental arsenic is toxic, as are many of its inorganic compounds; however some of its organic compounds can promote growth in chickens. The lethal dose of arsenic for a typical adult is 200 mg and can cause diarrhea, vomiting, colic, dehydration, and coma. Death from arsenic poisoning typically occurs within a day. Antimony is mildly toxic. Additionally, wine steeped in antimony containers can induce vomiting. When taken in large doses, antimony causes vomiting in a victim, who then appears to recover before dying several days later. Antimony attaches itself to certain enzymes and is difficult to dislodge. Stibine, or SbH3, is far more toxic than pure antimony. Bismuth itself is largely non-toxic, although consuming too much of it can damage the liver. Only one person has ever been reported to have died from bismuth poisoning. However, consumption of soluble bismuth salts can turn a person's gums black. Moscovium is too unstable to conduct any toxicity chemistry.	Physical sciences	Group 15	Chemistry
188044	https://en.wikipedia.org/wiki/Convertible	Convertible	A convertible or cabriolet () is a passenger car that can be driven with or without a roof in place. The methods of retracting and storing the roof vary across eras and manufacturers. A convertible car's design allows an open-air driving experience, with the ability to provide a roof when required. A potential drawback of convertibles is their reduced structural rigidity (requiring significant engineering and modification to counteract the side effects of almost completely removing a car's roof). The majority of convertible roofs are of a folding construction framework with the actual top made from cloth or other fabric. Other types of convertible roofs include retractable hardtops (often constructed from metal or plastic) and detachable hardtops (where a metal or plastic roof is manually removed and often stored in the trunk). Terminology Other terms for convertibles include cabriolet, cabrio, drop top, drophead coupé, open two-seater, open top, rag top, soft top, spider, and spyder, although companies use many of these terms interchangeably. Thus, nomenclatural consistency is rare. The term cabriolet originated from a carriage cabriolet: "a light, two-wheeled, one-horse carriage with a folding top, capable of seating two persons"; however, the term is also used to describe other convertibles. In the United Kingdom, the historical term for a two-door convertible is drophead coupé, and a four-door convertible was called an all-weather tourer. History Most of the early automobiles were open-air vehicles without any roof or sides. As car engines became more powerful by the end of the 19th century, folding textile or leather roofs (as had been used on victoria or landau carriages) began to appear on cars. Examples of early cars with roofs include the phaeton (a two-seat car with a temporary roof), the brougham or a coupé de ville, having an enclosed passenger compartment at the rear, while the driver sat in front either in the open, or the landaulet, where the driver has a fixed roof and the passenger compartment has a folding roof. Less expensive cars, such as the runabouts, sporting roadsters, or sturdy touring cars, remained either completely open air or were fitted with a rudimentary folding top and detachable clear side curtains. In the 1920s, when steel bodies began to be mass-produced, closed cars became available to the average buyer, and fully open cars began to disappear from the mainstream market. By the mid 1930s, the remaining small number of convertibles sold were high-priced luxury models. In 1939, Plymouth introduced the first mechanically operated convertible roof powered by two vacuum cylinders. Demand for convertibles increased as a result of American soldiers in France and the United Kingdom during World War II familiarizing themselves with small roadster cars, which were not available in the United States at that time. These roadsters included the MG Midget and Triumph Roadster. The convertible design was incorporated into the mass market unibody by Hudson in 1948. United States automakers manufactured a broad range of convertible models during the 1950s and 1960s – from economical compact-sized models such as the Rambler American and the Studebaker Lark, to the more expensive models, such as the Packard Caribbean, Oldsmobile 98, and Imperial by Chrysler. Automakers often included a convertible body style as an available body style in a model range. Convertibles in the U.S. market peaked in sales around 1965, and fell in popularity over the next five years. Optional air conditioning was gradually becoming more popular, and the availability of sunroofs and T-tops limited the appeal of the open body style. Noise, leaks, and repairs associated with fabric tops also contributed to issues that many customers had. The popularity of convertibles was reduced by the increased travel speeds on roads (resulting in more wind and noise for occupants) and the emergence of more comprehensive vehicle crash safety standards in the United States. The market share of convertibles fell to two or three percent of total sales and the U.S. automakers discontinued the body style from their lineups. American Motors stopped making convertibles after the 1968 model year, Chrysler after 1971, Ford after 1973, and most divisions of General Motors after 1975. Cadillac held out until 1976, when they made about 14,000. The last 200 had a red, white, and blue motif and a dashboard plaque. The very last was offered to the Smithsonian Institution, whose trustees turned it down as it was not at that moment a historic artifact, "Though it might well be in three generations ... or at the Tricentennial." After the last Cadillac Eldorado convertible was made in 1976, the only factory convertibles sold in the United States were imported. Making convertibles on the assembly line was both expensive and time-consuming, thus not worth the problems needed to sell the limited number of cars. Specialized coachbuilders were contracted to make dealer-available cars such as the Targa top versions of the AMC Concord and Eagle "Sundancer" as well as the Toyota Celica "Sunchaser" as specialty models. American Sunroof Company (ASC), which was responsible for popularizing the sunroof option for regular body styles, converted a Buick Riviera into a full convertible that compelled General Motors to market it as part of the 1982 Buick models. Chrysler Corporation also introduced a convertible body style in its 1982 lines that was based on the K-Car. These models were the LeBaron, produced under Chrysler, and the 400, manufactured under Dodge. Ford reintroduced a convertible Mustang for 1983, while American Motors Corporation (AMC) added a convertible version of the Renault Alliance in 1984. In 1989, Mazda released the first generation Mazda MX-5 (called "Miata" in North America), which has become the best-selling convertible with over 1 million cars sold. Also in 1989, Toyota released the Toyota Soarer Aerocabin, which uses an electrically operated retractable hardtop roof. A total of 500 were produced. Models dedicated to the convertible body style include the Mazda MX-5, Porsche Boxster, and Opel Cascada. Roof types Textile A "soft top" is made from a flexible textile material: Early convertibles used cotton canvas woven so tightly that it was waterproof. Automakers had problems in securing raw materials to fulfill orders after World War II, including canvas in various shades for convertible tops, therefore limiting their manufacture. A cloth-based material has become more common in recent years. Other materials are also used in the convertible top. By 1955, the most popular materials were latex and butyl rubber fabrics that each accounted for around 35% of the convertible top's weight, with others included vinyl (12%), jute (8%), along with rayon and acrylic fibers (Orlon), amounting to about 1% each in the compositions. Polyvinyl chloride (PVC) material was used for many convertible tops. The material consists of two layers: a top layer made of PVC, which has a specific structure depending on the vehicle model, and a lower layer made of fabric (usually cotton). The collapsible textile roof section over an articulated folding frame may include linings such as a sound-deadening layer and/or an interior cosmetic lining, to hide the frame. The folded convertible mechanism with the top is called the stack. Designs that fold down to a lower stack height offer a smoother silhouette for the car with the top down while concealed side rails allow room for three passengers in the back seat such as on the 1967 Rambler Rebel convertible. Detachable hardtop Rigid removable hardtops, many of which can be stored in a car's trunk/boot, have been available at least since the 1950s. These usually provide greater weatherproofing, soundproofing, and durability compared to fabric-based tops; some are available with integrated rear-window defrosters and windscreens. Examples include the Ford Thunderbird (1st-generation and 11th-generation), Mercedes SL (2nd-generation and 3rd-generation), Porsche Boxster, Jeep Wrangler, Ford Mustang Cobra (1995 Only), and Mazda MX-5. During the 1950s and 1960s, detachable hard-material roofs were offered for various convertible sports cars and roadsters, including the 1955–1957 Ford Thunderbird and Chevrolet Corvette, as well as the 1963–1971 Mercedes-Benz W113 series of two-seaters. Because the convertible top mechanism is itself expensive, the hard roof was customarily offered as an additional, extra-cost option. On early Thunderbirds (and Corvettes through 1967), buyers could choose between a detachable hardtop and a folding canvas top at no additional cost, but paid extra for both. The metal-framed "Carson top" was a popular addition for the 1930s Ford convertibles or roadsters because it turned these models into an almost instant hardtop. The design mimicked a convertible top, but lacking the bulky folding mechanisms enabled the removable hardtop to have a much lower and more rakish profile. Improvements in canvas tops have rendered the detachable hard roof less common in part because the top cannot be stored inside the vehicle when not in use, requiring a garage or other storage facility. Some open cars continue to offer it as an option. For example, the Mazda MX-5 has an accessory hardtop, which is compulsory for some auto racing series. Retractable hardtop A retractable hardtop — also known as "coupé convertible" or "coupé cabriolet" — is a car with an automatically operated, self-storing hardtop (as opposed to the textile-based roof used by traditional convertibles). The benefits of improved climate control and security are traded off against increased mechanical complexity, cost, weight, and often reduced luggage capacity. Other design features Tonneau cover Folding textile convertible tops often fail to completely hide their internal mechanism or can expose their vulnerable underside to sun exposure and fading. A tonneau cover provides a solution. Rear window Rear windows are often part of the roof assembly. Traditionally, the rear window in a soft-top was made from plastic; however, more recently some convertibles have used glass for the rear window. Windblocker A windblocker or wind deflector minimizes noise and rushing air reaching the occupants. According to the engineer responsible for the 2008 Chrysler Sebring, its windblocker reduces wind noise by approximately 11 to 12 dB. Several convertibles are available with a heating duct to the neck area of the seat, which is often called an "Air Scarf". Examples of cars with this feature include Mercedes-Benz SLK-Class, Mercedes-Benz SL-Class, and Audi A5/S5. Safety Modern safety features specifically for convertibles include: rollover protection structures (ROPS) with pyrotechnically charged roll hoops hidden behind the rear seats that deploy under rollover conditions heated rear window (for improved visibility) boron steel-reinforced A-pillars safety cage construction – a horseshoe-like structure around the passenger compartment door-mounted side-impact airbag which inflates upward (instead of downward like the typical curtain airbag) to provide head protection even with an open window Variations Convertibles have offered numerous iterations that fall between the first mechanically simple fabric tops to complex retractable roofs made from hard materials: Roadster: A roadster (also called spider or spyder) is an open two-seat car with emphasis on sporting appearance or character. Initially, an American term for a two-seat car with no weather protection, usage has spread internationally and has evolved to include two-seat convertibles. Cabrio coach: A cabrio coach (also called semi-convertible) has a retractable textile roof, similar to a traditional convertible. The difference is that a convertible often has the B-pillar, C-pillar and other bodywork removed. However, the cabrio-coach retains all bodywork to the top of the door frames and just replaces the roof skin with a retractable fabric panel. An advantage of a cabrio coach is that retaining more of the car's original structure means that structural rigidity is higher (or the vehicle weight is lower) than traditional cabriolets. An example of the cabrio coach is the 2003-10 C3 Pluriel, which has a roof with five possible configurations. Fixed-profile: In contrast to convertibles where the entire bodywork above the beltline (doors, roof, side pillars, side bodywork) is replaced with a folding or retractable roof, the fixed profile convertible retains portions of fixed bodywork including the doors, side pillars, and side elements of the roof — while a center fabric portion slides back and accordions at the rear. As an example, Citroën's 1948 Citroën 2CV featured rigid bodysides and two doors on each side, along with a sunroof that rolled back on itself and extended to the rear bumper in place of a separate trunk lid. Other fixed-profile convertibles include the 1957 Autobianchi Bianchina Trasformabile, 1957 Vespa 400, 1950 Nash Rambler Landau Convertible Coupe, the Nissan Figaro (1991), the Jaguar XJ-SC (1983) as well as the 1957 Fiat 500 and its 2007 Fiat 500 successor. The 1984 Heuliez-designed Citroën Visa Décapotable used elements of a fixed-profile convertible. Four-door: Most convertibles have two doors. However, four-door convertibles have been mass-produced. Examples include the 1940-41 Cadillac Series 62, 1931 Chrysler Imperial Dual Cowl Phaeton and 1961-67 Lincoln Continental. Current production four-door convertibles include the Jeep Wrangler Unlimited. Peugeot presented a concept four-door retractable hardtop convertible, the Peugeot 407 Macarena in 2006. Produced by French coachbuilding specialist Heuliez, the Macarena's top can be folded in 60 seconds, with a steel reinforcing beam behind the front seats incorporating LCD screens for the rear passengers into the crossmember. Off-road: Several off-road vehicles have been produced with removable soft tops. Examples include the Jeep Wrangler, Suzuki Vitara, Suzuki Jimny, Ford Bronco, Land Rover Defender, Mercedes-Benz G-Class as well as early models of the Toyota Land Cruiser and Land Rover Defender. Typically, the soft tops attach to the roll cage or to the installation points on the vehicle's body. Landaulet: A landaulet (also known as landaulette) is where the rear passengers are covered by a convertible top. Often the driver is separated from the rear passengers with a partition, as per a limousine. In the second half of the 20th century, landaulets were used by public figures (such as heads of state) in formal processions. They are now rarely used, for fear of terrorist attacks. Victoria-Cabriolet: reminiscent of the victoria carriage style, a three-position convertible. No rear side windows and equipped with a soft top that can be raised partway, leaving the area above the front seats folded back. This body style had a short period of popularity, mainly in the 1930s. Other names include Cabriolet/Coupé Milord (or just Milord), Calash (from Calèche), Folding Head DHC, three-position Drop-head Coupé, or Cabriolet toit de 3 positions. Gallery Open car and roadster Convertibles Retractable hardtop	Technology	Motorized road transport	null
188086	https://en.wikipedia.org/wiki/USS%20Arizona	USS Arizona	USS Arizona was a standard-type battleship built for the United States Navy in the mid-1910s. Named in honor of the 48th state, she was the second and last ship in the . After being commissioned in 1916, Arizona remained stateside during World War I but escorted President Woodrow Wilson to the subsequent Paris Peace Conference. The ship was deployed abroad again in 1919 to represent American interests during the Greco-Turkish War. Two years later, she was transferred to the Pacific Fleet, under which the ship would remain for the rest of her career. The 1920s and 1930s saw Arizona regularly deployed for training exercises, including the annual Fleet Problems, excluding a comprehensive modernization between 1929 and 1931. The ship supported relief efforts in the wake of a 1933 earthquake near Long Beach, California, and was later filmed for a role in the 1934 James Cagney film Here Comes the Navy before budget cuts led to significant periods in port from 1936 to 1938. In April 1940, the Pacific Fleet's home port was moved from California to Pearl Harbor, Hawaii, as a deterrent to Japanese imperialism. On 7 December 1941, the Japanese attacked Pearl Harbor, and Arizona was hit by several air-dropped armor-piercing bombs. One detonated an explosive-filled magazine, sinking the battleship and killing 1,177 of its officers and crewmen. Unlike many of the other ships attacked that day, Arizona was so irreparably damaged that it was not repaired for service in World War II. The shipwreck still lies at the bottom of Pearl Harbor beneath the USS Arizona Memorial. Dedicated to all those who died during the attack, the memorial is built across the ship's remains. Description The Pennsylvania-class ships were significantly larger than their predecessors, the . Arizona had an overall length of , a beam of (at the waterline), and a draft of at deep load. This was longer than the older ships. She displaced at standard and at deep load, over more than the older ships. The ship had a metacentric height of at deep load. Her crew numbered 56 officers and 1,031 enlisted men as built. The ship had four direct-drive Parsons steam turbine sets, each of which drove a propeller in diameter using steam provided by twelve Babcock & Wilcox boilers. The turbines were designed to produce a total of , but achieved only during Arizonas sea trials, when she met her designed speed of . However, she did manage to reach during a full-power trial in September 1924. She was designed to carry enough fuel oil to steam at a speed of for with a clean bottom. She had four turbo generators. Arizona carried twelve 45-caliber 14-inch guns in triple gun turrets. The turrets were numbered from I to IV from front to rear. Defense against torpedo boats was provided by twenty-two 51-caliber guns mounted in individual casemates in the sides of the ship's hull. Positioned as they were they proved vulnerable to sea spray and could not be worked in heavy seas. The ship mounted four 50-caliber guns for anti-aircraft defense, although only two were fitted when completed. The other pair was added shortly afterwards on top of Turret III. Arizona also mounted two torpedo tubes underwater, one on each broadside, and carried 24 torpedoes for them. The Pennsylvania-class design continued the all-or-nothing principle of armoring only the most important areas of the ship begun in the Nevada class. The waterline armor belt of Krupp armor measured thick and covered only the ship's machinery spaces and magazines. It had a total height of , of which was below the waterline; beginning below the waterline, the belt tapered to its minimum thickness of . The transverse bulkheads at each end of the ship ranged from 13 to 8 inches in thickness. The faces of the gun turrets were thick while the sides were thick and the turret roofs were protected by of armor. The armor of the barbettes was thick. The conning tower was protected by of armor and had a roof eight inches thick. The main armor deck was three plates thick with a total thickness of 3 inches; over the steering gear the armor increased to in two plates. Beneath it was the splinter deck that ranged from in thickness. The boiler uptakes were protected by a conical mantlet that ranged from in thickness. A three-inch torpedo bulkhead was placed inboard from the ship's side and the ship was provided with a complete double bottom. Testing in mid-1914 revealed that this system could withstand of TNT. Construction and trials The keel of battleship number 39 (hull number: BB-39) was laid on the morning of 16 March 1914 with Assistant Secretary of the Navy Franklin Delano Roosevelt in attendance. The builders intended to set a world-record ten months between the ship's keel-laying and launch, for what The New York Times declared would be "the world's biggest and most powerful, both offensively and defensively, superdreadnought ever constructed," but the ship was only a little over half complete a year later. She was launched on 19 June 1915, making it about fifteen months from keel-laying to launch. In the meantime, the ship was named after the newest state in the union by Secretary of the Navy Josephus Daniels. The New York Times estimated that 75,000 people attended the launch, including John Purroy Mitchel, the mayor of New York City, George W. P. Hunt, the governor of Arizona, and many high-ranking military officials. Several warships were also nearby, including many of the new dreadnoughts which had already entered service (, , , , , and ). Esther Ross, the daughter of W. W. Ross of Prescott, Arizona, was given the honors of ship sponsor and christening. To acknowledge a ban on alcohol recently passed by the state legislature, the state's governor decided that two bottles would be used: one full of sparkling wine from Ohio, and another filled with water from the Roosevelt Dam. After the launch, Arizona was towed to the Brooklyn Navy Yard for fitting-out. Arizona was commissioned into the Navy on 17 October 1916 with Captain John McDonald in command. She departed New York on 10 November 1916 after the crew had cleaned the ship and the propulsion system had been tested at the dock. After declinating the ship's magnetic compasses, the ship sailed south for her shakedown cruise. Outside Guantanamo Bay, a stripped turbine on 7 December forced the navy to order Arizona back to New York for repairs, although she was able to enter Chesapeake Bay to test her main and secondary gun batteries on 19–20 December. The turbine could not be repaired inside the ship, so the yard workers had to cut holes in the upper decks to lift the damaged casing out. It was reinstalled after almost four months of repairs at the naval yard. World War I Arizona left the yard on 3 April 1917, and three days later, the United States declared war on Germany. Assigned to Battleship Division 8 operating out of the York River, Arizona was employed only as a gunnery training ship for the crewmen on armed merchant vessels crossing the Atlantic in convoys. Shortly after the war began, eight of her 5-inch guns (the four guns farthest forward and the sternmost four guns) were removed to equip merchant ships. When the ship sailed near the wreck of the old San Marcos (ex-Texas), the wreck was sometimes used as a target for the 14-inch guns. Arizona rarely ventured into the ocean for fear of U-boats, and when she did, it was only in the company of other battleships and escort ships. Four coal-fired American dreadnoughts (it was easier to obtain coal than oil in the United Kingdom) were eventually sent across the Atlantic in December 1917 as Battleship Division Nine, but Arizona was not among them. Life for Arizonas crew was not all training, as the race-boat team from Arizona was able to win the Battenberg Cup in July 1918 by beating the team from by three lengths over the three-mile course. The fighting ended on 11 November 1918 with an armistice. A week later, the ship left the United States for the United Kingdom, arriving on 30 November 1918. After two weeks berthed at Portland Harbor in Dorset, Arizona sailed for France. On 13 December 1918, Arizona joined nine battleships and twenty-eight destroyers escorting President Woodrow Wilson on the ocean liner into Brest for one day on Wilson's journey to the Paris Peace Conference. The ten battleships departed France the next day, taking less than two weeks to cross the Atlantic, and arrived in New York on 26 December to parades, celebrations, and a full naval review by Secretary Daniels. Arizona was the first in line and rendered a nineteen-gun salute to Daniels. Along with many of the other members of the recently returned fleet, she was anchored off New York City for the next several weeks and open to the public. Post-war and the 1920s Arizona sailed from New York for Hampton Roads on 22 January 1919; she continued south to Guantanamo Bay on 4 February and arrived on four days later. The time in Caribbean waters was mostly used in training for battles and fleet maneuvering, although it included a liberty visit to Port of Spain. In April, Arizonas crew won the Battenberg Cup rowing competition for the second straight year before the ship was deployed to France once again to escort President Wilson back to the United States. While the ship was awaiting Wilson's departure, she was redeployed to Smyrna (now İzmir) in Turkey in response to tensions between Greece and Italy over the awarding of Smyrna to Greece in the Paris Peace Treaty. The Greek and Italian governments had each deployed a major warship to the area (Georgios Averof and Duilio, respectively) to enforce their interests. Shortly after Arizona arrived, Greek ground forces arrived in transports and were off-loaded in the port. The resultant chaos in the city caused many American citizens in the area to seek shelter on board Arizona. When the crisis abated, Arizona was ordered to Constantinople (now Istanbul) before she sailed for home on 15 June. She put into the New York Navy Yard on 30 June for an overhaul, where six 5-inch guns were removed and the fire control system was modernized. Work was completed in January 1920 and the battleship sailed south to Guantanamo Bay for crew training. During this time, Arizona was fitted with a flying-off platform similar to the one given to Texas in March 1919. In April, Arizona lost the Battenberg Cup to Nevada, and in June she was present for the Naval Academy's graduation ceremonies. In August she became the flagship of Battleship Division Seven, although it was only later in 1920 that the battleship was refitted to be an admiral's flagship. In company with six battleships and eighteen destroyers, Arizona was sent south again to transit the Panama Canal in January 1921. After meeting up with the Pacific Fleet, Arizona continued on to Peru for a week before the two fleets combined to practice battle maneuvers. After a short return to the Atlantic, which included an overhaul in New York, Arizona, under the command of Jehu V. Chase, returned to Peru in the summer before she began operating from her new home port of San Pedro, California, part of Los Angeles, where she was based until 1940. For the rest of the 1920s, Arizonas service consisted of routine training exercises. Naval historian Paul Stillwell remarked that "the Pacific years included a great deal of sameness and repetition", and his chronology of the ship's movements is filled with phrases like "torpedo-defense practice", "battle-practice rehearsal", "gunnery practice", "en route to…", and "anchored at…". A recurring theme in these years were the annual Fleet Problems, which began in 1923 and simulated large fleet actions by having most of the active fleet face off against each other. The first two simulated an attack on the Panama Canal from the west, while in 1925 they attempted to defend the Hawaiian Islands. Other 1920s Fleet Problems included the Caribbean, near Central America, the West Indies, and Hawaii. On 27 July 1923 the ship, under command of John Y.R. Blakely, joined President Warren G. Harding's naval review in Seattle. Harding died just one week later, and Arizona joined the Pacific Fleet to fire a salute in his honor on 3 August. Sometime in early March 1924 a prostitute named Madeline Blair stowed away aboard Arizona, trading sex for a free voyage to San Pedro until she was discovered on 12 April while the ship was anchored in Balboa, Panama. She was sent back to New York City and Captain Percy Olmstead later convened courts-martial for 23 sailors once the ship began her refit in the Bremerton Navy Yard, which imposed sentences of up to 10 years imprisonment. Admiral Henry A. Wiley, commander of the Battle Fleet, issued a letter of reprimand to all officers of the ship, including future Admiral and Chief of Naval Operations Arleigh Burke, then an ensign. Admiral William V. Pratt, then in command of the division to which Arizona was assigned, thought the penalties excessive, and he ordered the reprimands stricken from the officer's records when he became Chief of Naval Operations in 1930. Modernization Four months after Fleet Problem IX in January 1929, Arizona was modernized at the Norfolk Navy Yard. New tripod masts, surmounted by three-tiered fire-control directors for the main and secondary armament, replaced the old hyperboloid cage masts; the number of five-inch guns was reduced to 12 and the guns re-positioned one deck higher, and eight 25-caliber five-inch anti-aircraft guns replaced the three-inch guns with which she had been originally equipped. These changes increased her crew to 92 officers and 1,639 enlisted men. The ship's main gun turrets were modified to increase the maximum elevation of their guns to 30°. The compressed-air catapult on the quarterdeck was replaced by one that used black powder. Her deck armor was increased by the addition of a thickness of Special Treatment Steel, and the ship was bulged to protect her from torpedoes. An additional bulkhead was added to the sides of the boiler rooms for the same purpose. At the same stroke, her own outfit of two submerged torpedo tubes was removed during this refit in light of a new appreciation that anticipated battleship engagement ranges made their future use improbable. This alteration also permitted the large transverse flat in which the tubes had been situated to be subdivided to reduce risks of flooding in action. Arizonas machinery was almost entirely replaced; her high-pressure turbines were replaced by more powerful geared turbines from the canceled battleship , and six new boilers replaced her originals. Their additional power offset the ship's increased displacement as demonstrated during her sea trials; Arizona made with at a displacement of . 1930s On 19 March 1931, even before Arizona was put through post-modernization sea trials, she hosted President Herbert Hoover for a brief vacation in the Caribbean. The President visited Puerto Rico and the Virgin Islands. Returning on 29 March, Arizona conducted her sea trials at Rockland, Maine, and had another catapult fitted on the top of Turret III, before she was transferred to the West Coast in August with her sister Pennsylvania. In February 1932, the ship participated in Grand Joint Exercise No. 4 in which carrier aircraft successfully attacked Pearl Harbor on Sunday morning, 7 February. After returning to the West Coast from Fleet Problem XIV in 1933, the ship was anchored in San Pedro when an earthquake struck nearby Long Beach, California, on 10 March. Sailors from the ship joined the relief efforts, providing food, treating the injured, and providing security from looters. In early 1934, the ship and her crew were filmed for the James Cagney/Warner Brothers film Here Comes the Navy, which made extensive use of exterior footage as well as on-board location shots. In the early morning of 26 July, Arizona collided with a fishing trawler, Umatilla, that was under tow by another trawler off Cape Flattery. Two men aboard Umatilla were killed in the collision and the Navy convened a Court of Inquiry to investigate the incident. The court recommended that the ship's captain, Captain MacGillivray Milne, be court-martialed. This took place at Guantanamo Bay Naval Base, Cuba, while the ship was participating in that year's Fleet Problem off the East Coast. Milne was judged guilty and replaced several months later by Captain George Baum after the ship returned to the West Coast. In the meantime, Rear Admiral Samuel W. Bryant assumed command of Battleship Division Two on 4 September, with Arizona as his flagship. Rear Admiral George T. Pettengill relieved Bryant on 4 March 1935 and the ship participated in Fleet Problem XVI two months later. Arizona made a port visit to Balboa in May 1936 during Fleet Problem XVII. On 8 June, Captain George A. Alexander relieved Baum as captain, and, 15 days later, Rear Admiral Claude C. Bloch relieved Pettengill. During gunnery practice on 24 July, the combustion gases from one gun of Turret II entered the gun turret, burning one crewman. The turret's sprinkling system was turned on to prevent any powder explosion, but the released water leaked into the turret's electrical switchboard and started a small fire that was easily put out. Due to the navy's limited budget, the ship spent most of this period in port as a fuel-saving measure. In Fiscal Year 1936–37, the ship was anchored for 267 days; the following year it was in port for 255 days. The ship spent the rest of her career based on the West Coast or in Hawaii. On 2 January 1937, Rear Admiral John Greenslade assumed command of Battleship Division Two from Bloch and transferred his flag to the battleship on 13 April. Rear Admiral Manley H. Simons, commander of Battleship Division One, transferred his flag to Arizona on 7 August. He was relieved by Rear Admiral Adolphus E. Watson on 8 November. Captain Alfred Winsor Brown relieved Baum on 11 December. The ship participated in Fleet Problem XIX off Hawaii in April–May 1938. Captain Brown died in his sleep on 7 September and Captain Isaac C. Kidd assumed command of the ship on 17 September 1938. That same day, Rear Admiral Chester Nimitz assumed command of Battleship Division One. Nimitz was relieved on 27 May 1939 by Rear Admiral Russell Willson. Captain Harold C. Train assumed command of the ship on 5 February 1940. Arizonas last fleet problem was off Hawaii in April–May 1940. At its conclusion, the United States Pacific Fleet was retained in Hawaiian waters, based at Pearl Harbor, to deter the Japanese. She was overhauled at the Puget Sound Navy Yard, Bremerton, Washington, from October 1940 to January 1941. During this refit, the foundation for a search radar was added atop her foremast, her anti-aircraft directors were upgraded and a platform for four water-cooled caliber M2 Browning machine guns was installed at the very top of the mainmast. Her last flag change-of-command occurred on 23 January 1941, when Willson was relieved by Isaac Kidd, by that time a rear admiral. Captain Franklin Van Valkenburgh relieved Train on 5 February 1941. On 22 October 1941, during an exercise taking place in heavy fog, the ship was hit in the bow by the . Arizona had been scheduled to depart for Bremerton Navy Yard in November to undergo an overhaul. The accident instead required her to be dry-docked at Pearl Harbor for repairs to the collision damage. As a result, she remained in Hawaii. The ship's last sortie was a night-firing exercise on the night of 4 December as part of Battleship Division One, alongside Nevada and Oklahoma. All three ships moored at quays along Ford Island on the following day. On 6 December, the repair ship came alongside to assist the ship's crew with minor repairs. Attack on Pearl Harbor Shortly before 08:00 local time on 7 December 1941, Japanese aircraft from six aircraft carriers struck the Pacific Fleet as it lay in port at Pearl Harbor, and wreaked devastation on the warships and installations defending Hawaii. On board Arizona, the ship's air raid alarm went off at about 07:55, and the ship went to general quarters soon after. Shortly after 08:00, ten Nakajima B5N2 "Kate" torpedo bombers, five each from the carriers and , attacked Arizona. All of the aircraft were carrying armor-piercing shells modified into bombs. Flying at an estimated altitude of , Kagas aircraft bombed Arizona from amidships to stern. Soon after, Hiryūs bombers hit the bow area. The aircraft scored four hits and three near-misses on and around Arizona. The near-miss off the port bow is thought to have caused observers to believe that the ship had been torpedoed, although no torpedo damage has been found. The stern-most bomb ricocheted off the face of Turret IV and penetrated the deck to detonate in the captain's pantry, causing a small fire. The next forward-most hit was near the port edge of the ship, abreast the mainmast, probably detonating in the area of the anti-torpedo bulkhead. The next bomb struck near the port rear 5-inch AA gun. Magazine explosion The last bomb hit at 08:06 in the vicinity of Turret II, likely penetrating the armored deck near the magazines located in the forward section of the ship. While not enough of the ship is intact to judge the exact location, its effects are indisputable: about seven seconds after the hit, the forward magazines detonated in a cataclysmic explosion, mostly venting through the sides of the ship and destroying much of the interior structure of the forward part of the ship. This caused the forward turrets and conning tower to collapse downward some and the foremast and funnel to collapse forward, effectively tearing the ship in two. The explosion touched off fierce fires that burned for two days; debris showered down on Ford Island in the vicinity. The blast from this explosion also put out fires on the repair ship Vestal, which was moored alongside. The bombs and subsequent explosion killed 1,177 of the 1,512 crewmen on board at the time, approximately half of the lives lost during the attack. Two competing hypotheses have arisen about the cause of the explosion. The first is that the bomb detonated in or near the black-powder magazine used for the ship's saluting guns and catapult charges. This would have detonated first and then ignited the smokeless powder magazines which were used for the ship's main armament. A 1944 Navy Bureau of Ships report suggests that a hatch leading to the black powder magazine was left open, possibly with flammable materials stocked nearby. The Naval History and Heritage Command explained that black powder might have been stockpiled outside the armored magazine. The alternative explanation is that the bomb penetrated the armored decks and detonated directly inside one of the starboard magazines for the main armament, but smokeless powder is relatively difficult to detonate. Thus the 14-inch powder bags required a black powder pad to quickly ignite the powder. The time elapsed from the bomb hit to the magazine explosion was shorter than experience suggested burning smokeless powder required to explode. It seems unlikely that a definitive answer to this question will ever be found, as the surviving physical evidence is insufficient to determine the cause of the magazine explosion. Awards and recognition After the attack, several sailors received medals for their conduct and actions under fire. Lieutenant Commander Samuel G. Fuqua, the ship's damage control officer, earned the Medal of Honor for his cool-headedness while quelling fires and getting survivors off the wrecked battleship. Posthumous awards of the Medal of Honor also went to two high-ranking officers who were on board the battleship when it was destroyed: Rear Admiral Kidd, the first flag officer killed in the Pacific war, and Captain Van Valkenburgh, who reached the bridge and was attempting to defend his ship when the bomb that hit the onboard ammunition magazines destroyed it. Arizona was awarded one battle star for her service in World War II. Salvage and memorial Arizona was placed "in ordinary" (declared to be temporarily out of service) at Pearl Harbor on 29 December, and was stricken from the Naval Vessel Register on 1 December 1942. She was so badly damaged by the magazine explosion that she was not thought fit for service even if she could be salvaged, unlike many of the other sunken ships nearby. Her surviving superstructure was scrapped in 1942, and her main armament was salvaged over the next year and a half. The aft main gun turrets were removed and reinstalled as United States Army Coast Artillery Corps Battery Arizona at Kahe Point on the west coast of Oahu and Battery Pennsylvania on the Mokapu Peninsula, covering Kaneohe Bay at what is now Marine Corps Base Hawaii. Battery Pennsylvania fired its guns for the first and last time on V-J Day in August 1945 while training, while the nearby Battery Arizona was never completed. Both forward turrets were left in place, although the guns from Turret II were salvaged and later installed on Nevada in the fall of 1944 after having been straightened and relined. Nevada later fired these same guns against the Japanese islands of Okinawa and Iwo Jima. Arizona memorials It is commonly—but incorrectly—believed that Arizona remains perpetually in commission, like the . Arizona is under the control of the National Park Service, but the US Navy still retains the title. Arizona retains the right, in perpetuity, to fly the United States flag as if she were an active, commissioned naval vessel. The wreck of Arizona remains at Pearl Harbor to commemorate the men of her crew lost that December morning in 1941. On 7 March 1950, Admiral Arthur W. Radford, commander in chief of the Pacific Fleet at that time, instituted the raising of colors over her remains. Legislation during the administrations of presidents Dwight D. Eisenhower and John F. Kennedy resulted in the designation of the wreck as a national shrine in 1962. A memorial was built across the ship's sunken remains, including a shrine room listing the names of the lost crew members on a marble wall. The national memorial was administratively listed on the National Register of Historic Places on 15 October 1966. The ship herself was designated a National Historic Landmark on 5 May 1989. Upon their death, survivors of the attack were able to have their ashes placed within the ship among their fallen comrades. Veterans who served aboard the ship at other times had the choice of scattering their ashes in the water above the ship. The last survivor of Arizona, Lou Conter, died in April 2024 at the age of 102. While the superstructure and two of the four main gun turrets were removed, the barbette of one of the turrets remains visible above the water. Since her sinking, oil still leaks from the hull, with more than escaping into the harbor per day. In 2004, the US Navy and the National Park Service oversaw a comprehensive computerized mapping of the hull, being careful to honor its role as a war grave. The navy considered non-intrusive means of abating the continued leakage of oil to avoid the further environmental degradation of the harbor. One of the original Arizona bells now hangs in the University of Arizona Student Union Memorial Center bell tower. The bell was rung after every home football victory over any team except other Arizona schools. , the bell is no longer rung due to the risk of damaging it. A gun, mast, and anchor from Arizona are in Wesley Bolin Memorial Plaza just east of the Arizona State Capitol complex in downtown Phoenix, Arizona. The gun's plaque states that it was not on the ship during the Pearl Harbor attack, but was being relined for mounting on the battleship Nevada. It is paired with a gun from the battleship to represent the start and end of the Pacific War for the United States. Other artifacts from the ship, such as items from the ship's silver service, are on permanent exhibit in the Arizona State Capitol Museum. Every two years the Navy awards "The USS Arizona Memorial Trophy" to a ship that has achieved the highest combat readiness in Strike warfare, Surface Fire Support and Anti-Surface warfare, as determined by the Chief of Naval Operations. The bronze trophy on a black marble base was provided to the Navy by the citizens of the state of Arizona on 7 December 1987.	Technology	Specific seacraft	null
188213	https://en.wikipedia.org/wiki/Yahoo	Yahoo	Yahoo (, styled yahoo! in its logo) is an American web services portal. The web portal provides search engine Yahoo Search and related services including My Yahoo, Yahoo Mail, Yahoo News, Yahoo Finance, Yahoo Sports and its advertising platform, Yahoo Native. It is operated by the namesake company Yahoo! Inc., which is 90% owned by Apollo Global Management and 10% by Verizon. Yahoo was established by Jerry Yang and David Filo in January 1994 and was one of the pioneers of the early Internet era in the 1990s. However, its use declined in the 2010s as some of its services were discontinued, and it lost market share to Facebook and Google. Etymology The word "yahoo" is a backronym for "Yet Another Hierarchically Organized Oracle" or "Yet Another Hierarchical Officious Oracle". The term "hierarchical" described how the Yahoo database was arranged in layers of subcategories. The term "oracle" was intended to mean "source of truth and wisdom", and the term "officious", rather than being related to the word's normal meaning, described the many office workers who would use the Yahoo database while surfing from work. However, founders Filo and Yang insist they mainly selected the name because they liked the slang definition of a "yahoo" (used by college students in David Filo's native Louisiana in the late 1980s and early 1990s to refer to an unsophisticated, rural Southerner): "rude, unsophisticated, uncouth." This meaning derives from the Yahoo race of fictional beings from Gulliver's Travels. History Founding In January 1994, Jerry Yang and David Filo were electrical engineering graduate students at Stanford University, when they created a website named "Jerry and David's guide to the World Wide Web". The site was a human-edited web directory, organized in a hierarchy, as opposed to a searchable index of pages. In March 1994, "Jerry and David's Guide to the World Wide Web" was renamed "Yahoo!" and became known as the Yahoo Directory. The "yahoo.com" domain was registered on January 18, 1995. Yahoo was incorporated on March 2, 1995. In 1995, a search engine function, called Yahoo Search, was introduced. This allowed users to search Yahoo Directory. Yahoo soon became the first popular online directory and search engine on the World Wide Web. Expansion Yahoo grew rapidly throughout the 1990s. Yahoo became a public company via an initial public offering in April 1996 and its stock price rose 600% within two years. Like many search engines and web directories, Yahoo added a web portal, putting it in competition with services including Excite, Lycos, and America Online. By 1998, Yahoo was the most popular starting point for web users, and the human-edited Yahoo Directory the most popular search engine, receiving 95 million page views per day, triple that of rival Excite. It also made many high-profile acquisitions. Yahoo began offering free e-mail from October 1997 after the acquisition of RocketMail, which was then renamed to Yahoo Mail. In 1998, Yahoo replaced AltaVista as the crawler-based search engine underlying the Directory with Inktomi. Yahoo's two biggest acquisitions were made in 1999: Geocities for $3.6 billion and Broadcast.com for $5.7 billion. Its stock price skyrocketed during the dot-com bubble, closing at an all-time high of $118.75/share on January 3, 2000. However, after the dot-com bubble burst, it reached a post-bubble low of $8.11 on September 26, 2001. Yahoo began using Google for search in June 2000. Over the next four years, it developed its own search technologies, which it began using in 2004 partly using technology from its $280 million acquisition of Inktomi in 2002. In response to Google's Gmail, Yahoo began to offer unlimited email storage in 2007. In 2008, the company laid off hundreds of people as it struggled from competition. In February 2008, Microsoft made an unsolicited bid to acquire Yahoo for $44.6 billion. Yahoo rejected the bid, claiming that it "substantially undervalues" the company and was not in the interest of its shareholders. Although Microsoft increased its bid to $47 billion, Yahoo insisted on another 10%+ increase to the offer and Microsoft cancelled the offer in May 2008. Carol Bartz, who had no previous experience in Internet advertising, replaced Yang as CEO in January 2009. In September 2011, after failing to meet targets, she was fired by chairman Roy J. Bostock; CFO Tim Morse was named as Interim CEO of the company. In April 2012, after the appointment of Scott Thompson as CEO, several key executives resigned, including chief product officer Blake Irving. On April 4, 2012, Yahoo announced 2,000 layoffs, or about 14% of its 14,100 workers by the end of year, expected to save around $375 million annually. In an email sent to employees in April 2012, Thompson reiterated his view that customers should come first at Yahoo. He also completely reorganized the company. On May 13, 2012, Thompson was fired and was replaced on an interim basis by Ross Levinsohn, recently appointed head of Yahoo's new Media group. Several associates of Third Point Management, including Daniel S. Loeb were nominated to the board of directors. Thompson's total compensation for his 130-day tenure with Yahoo was at least $7.3 million. On July 15, 2012, Marissa Mayer was appointed president and CEO of Yahoo, effective July 17, 2012. In June 2013, Yahoo acquired blogging site Tumblr for $1.1 billion in cash, with Tumblr's CEO and founder David Karp continuing to run the site. In July 2013, Yahoo announced plans to open an office in San Francisco. On August 2, 2013, Yahoo acquired Rockmelt; its staff was retained, but all of its existing products were terminated. Data collated by comScore during July 2013 revealed that, during the month, more people in the U.S. visited Yahoo websites than Google; the first time that Yahoo outperformed Google since 2011. The data did not count mobile usage, nor Tumblr. Mayer also hired Katie Couric to be the anchor of a new online news operation and started an online food magazine. However, by January 2014, doubts about Mayer's progress emerged when Mayer fired her own first major hire, Henrique de Castro. On December 12, 2014, Yahoo acquired video advertising provider BrightRoll for $583 million. On November 21, 2014, Yahoo acquired Cooliris. In August 2023, it was announced Yahoo had acquired the San Francisco-headquartered social investing platform, Commonstock. In April 2024, it was announced Yahoo had acquired the AI-driven news aggregator app, Artifact. Decline, security breaches, and sale By December 2015, Mayer was criticized as performance declined. Mayer was ranked as the least likable CEO in tech. On February 2, 2016, Mayer announced layoffs amounting to 15% of the Yahoo workforce. On July 25, 2016, Verizon Communications announced the acquisition of Yahoo's core Internet business for $4.83 billion. The deal excluded Yahoo's 15% stake in Alibaba Group and 35.5% stake in Yahoo Japan. On February 21, 2017, as a result of the Yahoo data breaches, Verizon lowered its purchase price for Yahoo by $350 million and reached an agreement to share liabilities regarding the data breaches. On June 13, 2017, Verizon completed the acquisition of Yahoo and Marissa Mayer resigned. Yahoo, AOL, and HuffPost were to continue operating under their own names, under the umbrella of a new company, Oath Inc., later called Verizon Media. The parts of the original Yahoo! Inc. which were not purchased by Verizon Communications were renamed Altaba, which was later liquidated, making a final distribution in October 2020. In September 2021, investment funds managed by Apollo Global Management acquired 90% of Yahoo. In November 2021, Yahoo announced that it was ending operations in mainland China due to the increasingly challenging business and legal environment. Previously, the company discontinued China Yahoo Mail on August 20, 2013. In 2023, Yahoo announced that it would cut 20% of its workforce. The move followed mass layoffs from other tech giants including Google, Microsoft, Twitter, Inc, Meta, and Amazon. The company is set to lay off roughly 1,000 staff members of their 8,600 workers. Products and services For a list of all current and defunct services offered by Yahoo, see List of Yahoo-owned sites and services. Data breaches On September 22, 2016, Yahoo disclosed a data breach that occurred in late 2014, in which information associated with at least 500 million user accounts, one of the largest breaches reported to date. The United States indicted four men, including two employees of Russia's Federal Security Service (FSB), for their involvement in the hack. On December 14, 2016, the company revealed that another separate data breach had occurred in 2014, with hackers obtaining sensitive account information, including security questions, to at least one billion accounts. The company stated that hackers had utilized stolen internal software to forge HTTP cookies. On October 3, 2017, the company stated that all 3 billion of its user accounts were affected by the August 2013 theft. Criticism DMCA notice to whistleblower On November 30, 2009, Yahoo was criticized by the Electronic Frontier Foundation for sending a DMCA notice to whistleblower website "Cryptome" for publicly posting details, prices, and procedures on obtaining private information pertaining to Yahoo's subscribers. Censorship of private emails affiliated with Occupy Wall Street protests After some concerns over censorship of private emails regarding a website affiliated with Occupy Wall Street protests were raised, Yahoo responded with an apology and explained it as an accident. Partners and sponsorships On September 11, 2001, Yahoo announced its partnership with FIFA for the 2002 FIFA World Cup and 2006 FIFA World Cup tournaments. It was one of FIFA's 15 partners at the tournaments. The deal included co-branding the organization's websites. Yahoo sponsored the 2012 Sundance Film Festival. NBC Sports Group aligned with Yahoo Sports the same year with content and program offerings on mobile and desktop platforms. Yahoo announced television video partnerships in 2013 with Condé Nast, WWE, ABC NEWS, and CNBC. Yahoo entered into a 10-year collaboration in 2014, as a founding partner of Levi's Stadium, home of the San Francisco 49ers. The National Basketball Association partnered with Yahoo Sports to stream games, offer virtual and augmented-reality fan experiences, and in 2018 NBA League Pass. Yahoo Sportsbook launched in November 2019, a collaboration with BetMGM. BuzzFeed acquired HuffPost from Yahoo in November 2020, in a stock deal with Yahoo as a minority shareholder. The NFL partnered with Yahoo in 2020, to introduce a new "Watch Together" function on the Yahoo Sports app for interactive co-viewing through a synchronized livestream of local and primetime NFL games. The Paley Center for Media collaborated with Verizon Media to exclusively stream programs on Yahoo platforms beginning in 2020. Yahoo became the main sponsor for the Pramac Racing team and the first title sponsor for the 2021 ESport/MotoGP Championship season. Yahoo, the official partner for the September 2021 New York Fashion Week event also unveiled sponsorship for the Rebecca Minkoff collection via a NFT space. In September 2021, it was announced that Yahoo partnered with Shopify, connecting the e-commerce merchants on Yahoo Finance, AOL and elsewhere.	Technology	Portals/Platform sites	null
188288	https://en.wikipedia.org/wiki/Arteriosclerosis	Arteriosclerosis	Arteriosclerosis, literally meaning "hardening of the arteries", is an umbrella term for a vascular disorder characterized by abnormal thickening, hardening, and loss of elasticity of the walls of arteries; this process gradually restricts the blood flow to one's organs and tissues and can lead to severe health risks brought on by atherosclerosis, which is a specific form of arteriosclerosis caused by the buildup of fatty plaques, cholesterol, and some other substances in and on the artery walls (it can be brought on by smoking, a bad diet, or many genetic factors). Atherosclerosis is the primary cause of coronary artery disease (CAD) and stroke, with multiple genetic and environmental contributions. Genetic-epidemiologic studies have identified a long list of genetic and non-genetic risk factors for CAD. However, such studies indicate that family history is the most significant independent risk factor. Signs and symptoms The signs and symptoms of arteriosclerosis depend on the vessel affected by the disease. If affecting cerebral or ophthalmic vessels, as in cerebrovascular accidents or transient ischemic attacks, signs and symptoms may include sudden weakness, facial or lower limb numbness, confusion, difficulty understanding speech, and problems seeing. If affecting coronary vessels, as in coronary artery disease (including acute myocardial ischemia or a "heart attack"), signs and symptoms may include chest pain. Pathophysiology The lesions of arteriosclerosis begin as the intima (innermost layer of blood vessel wall) of an artery start to fill up with the deposition of cellular wastes. As these start to mature, they can take different forms of arteriosclerosis. All are linked through common features such as the stiffening of arterial vessels, thickening of arterial walls and the degenerative nature of the disease. Arteriolosclerosis, unlike atherosclerosis, is a sclerosis that only affects small arteries and arterioles, which carry nutrients and blood to the cells. Atherosclerosis is the narrowing of arteries from a buildup of plaque, usually made up of cholesterol, fatty substances, cellular waste products, calcium, and fibrin, inside the arteries. This affects large and medium-sized arteries; however, its positioning varies person to person. Monckeberg's arteriosclerosis or medial calcific sclerosis is seen mostly in the elderly, commonly in arteries of the extremities. Hyperplastic: Hyperplastic arteriosclerosis refers to the type of arteriosclerosis that affects large and medium-sized arteries. Hyaline type: Hyaline arteriosclerosis, also referred to as arterial hyalinosis and arteriolar hyalinosis, refers to lesions that are caused by the deposition of homogenous hyaline in the small arteries and arterioles. Diagnosis Diagnosis of an individual suspected of having arteriosclerosis can be based on a physical exam, blood test, EKG and the results of these tests (among other exams). Treatment Treatment is often in the form of preventive measures of prophylaxis. Medical therapy is often prescribed to help prevent arteriosclerosis for underlying conditions, such as medications for the treatment of high cholesterol (e.g., statins, cholesterol absorption inhibitors), medications to treat high blood pressure (e.g., ACE inhibitors, angiotensin II receptor blockers), and antiplatelet medications. Lifestyle changes are also advised, such as increasing exercise, stopping smoking, and moderating alcohol intake. There are a variety of types of surgery: Angioplasty and stent placement: A catheter is first inserted into the blocked or narrowed part of the artery, followed by a second one with a deflated balloon that is passed through the catheter into the narrowed area. The balloon is then inflated, pushing the deposits back against the arterial walls, and then a mesh tube is usually left behind to prevent the artery from retightening. Coronary artery bypass surgery: This surgery creates a new pathway for blood to flow to the heart. The surgeon attaches a healthy piece of vein to the coronary artery, just above and below the blockage to allow bypass. Endarterectomy: This is the general procedure for the surgical removal of plaque from the artery that has become narrowed or blocked. Thrombolytic therapy: This is a treatment used to break up masses of plaque inside the arteries via intravenous clot-dissolving medicine. Epidemiology In 2008, the US had an estimate of 16 million atherosclerotic heart disease and 5.8 million strokes. Cardiovascular diseases that were caused by arteriosclerosis also caused almost 812,000 deaths in 2008, more than any other cause, including cancer. About 1.2 million Americans are predicted to have a heart attack each year. History The diagnostics and clinical implications of this disease were not recognized until the 20th century. Many cases have been observed and recorded, and Jean Lobstein coined the term arteriosclerosis while he was analyzing the composition of calcified arterial lesions. The name "arteriosclerosis" is derived the Greek words ἀρτηρία (artēría, artery) and σκληρωτικός (sklērōtikós, hardened).	Biology and health sciences	Cardiovascular disease	Health
188325	https://en.wikipedia.org/wiki/Ureter	Ureter	The ureters are tubes composed of smooth muscle that transport urine from the kidneys to the urinary bladder. In an adult human, the ureters typically measure 20 to 30 centimeters in length and about 3 to 4 millimeters in diameter. They are lined with urothelial cells, a form of transitional epithelium, and feature an extra layer of smooth muscle in the lower third to aid in peristalsis. The ureters can be affected by a number of diseases, including urinary tract infections and kidney stone. is when a ureter is narrowed, due to for example chronic inflammation. Congenital abnormalities that affect the ureters can include the development of two ureters on the same side or abnormally placed ureters. Additionally, reflux of urine from the bladder back up the ureters is a condition commonly seen in children. The ureters have been identified for at least two thousand years, with the word "ureter" stemming from the stem relating to urinating and seen in written records since at least the time of Hippocrates. It is, however, only since the 1500s that the term "ureter" has been consistently used to refer to the modern structure, and only since the development of medical imaging in the 1900s that techniques such as X-ray, CT, and ultrasound have been able to view the ureters. The ureters are also seen from the inside using a flexible camera, called ureteroscopy, which was first described in 1964. Structure The ureters are tubular structures, approximately in adults, that pass from the pelvis of each kidney into the bladder. From the renal pelvis, they descend on top of the psoas major muscle to reach the brim of the pelvis. Here, they cross in front of the common iliac arteries. They then pass down along the sides of the pelvis and finally curve forward and enter the bladder from its left and right sides at the back of the bladder. The ureters are in diameter and surrounded by a layer of smooth muscle for near their ends just before they enter the bladder. The ureters enter the bladder from its back surface, traveling before opening into the bladder at an angle on its outer back surface at the slit-like ureteric orifices. This location is also called the vesicoureteric junction. In the contracted bladder, they are about apart and about the same distance from the internal urethral orifice; in the distended bladder, these measurements may be increased to about . A number of structures pass by, above, and around the ureters on their path down from the kidneys to the bladder. In its upper part, the ureter travels on the psoas major muscle and sits just behind the peritoneum. As it passes down the muscle, it travels over the genitofemoral nerve. The inferior vena cava and the abdominal aorta sit to the midline of the right and left ureters, respectively. In the lower part of the abdomen, the right ureter sits behind the lower mesentery and the terminal ileum, and the left ureter sits behind the jejunum and the sigmoid colon. As the ureters enter the pelvis, they are surrounded by connective tissue, and travel backward and outward, passing in front of the internal iliac arteries and internal iliac veins. They then travel inward and forward, crossing the umbilical, inferior vesical, and middle rectal arteries. From here, in males, they cross under the vas deferens and in front of the seminal vesicles to enter the bladder near the trigone. In females, the ureters pass behind the ovaries and then travel in the lower midline section of the broad ligament of the uterus. For a short part, the uterine arteries travel on top for a short () period. They then pass by the cervix, traveling inward towards the bladder. Blood and lymphatic supply The arteries which supply the ureter vary along its course. The upper third of the ureter, closest to the kidney, is supplied by the renal arteries. The middle part of the ureter is supplied by the common iliac arteries, direct branches from the abdominal aorta, and gonadal arteries; the gonadal arteries being the testicular artery in men and the ovarian artery in women. The lower third of the ureter, closest to the bladder, is supplied by branches from the internal iliac arteries, mainly the superior and inferior vesical arteries. The arterial supply can be variable, with arteries that contribute include the middle rectal artery, branches directly from the aorta, and, in women, the uterine and vaginal arteries. The arteries that supply the ureters end in a network of vessels within the adventitia of the ureters. There are many connections () between the arteries of the ureter, particularly in the adventitia, which means damage to a single vessel does not compromise the blood supply of the ureter. Venous drainage mostly parallels that of the arterial supply; that is, it begins as a network of smaller veins in the adventitia; with the renal veins draining the upper ureters, and the vesicular and gonadal veins draining the lower ureters. Lymphatic drainage depends on the position of lymphatic vessels in the ureter. Lymph collects in submucosal, intramuscular and adventitial lymphatic vessels. Those vessels closer to the kidney drain into renal collecting vessels, and from here into the lateral aortic nodes near the gonadal vessels. The middle part of the ureter drains into the right paracaval and interaortocaval nodes on the right, and the left paraaortic nodes on the left. In the lower ureter, lymph may drain into the common iliac lymph nodes, or lower down in the pelvis to the common, external, or internal iliac lymph nodes. Nerve supply The ureters are richly supplied by nerves that form a network () of nerves, the ureteric plexus that lies in the adventitia of the ureters. This plexus is formed from a number of nerve roots directly (T9–12, L1, and S2-4), as well as branches from other nerve plexuses and nerves; specifically, the upper third of the ureter receives nerve branches from the renal plexus and aortic plexus, the middle part receives branches from the upper hypogastric plexus and nerve, and the lower ureter receives branches from the lower hypogastric plexus and nerve. The plexus is in the adventitia. These nerves travel in individual bundles and along small blood vessels to form the ureteric plexus. Sensation supplied is sparse close to the kidneys and increases closer to the bladder. Sensation to the ureters is provided by nerves that come from T11 – L2 segments of the spinal cord. When pain is caused, for example by spasm of the ureters or by a stone, the pain may be referred to the dermatomes of T11 – L2, namely the back and sides of the abdomen, the scrotum (males) or labia majora (females) and upper part of the front of the thigh. Microanatomy The ureter is lined by urothelium, a type of transitional epithelium that is capable of responding to stretches in the ureters. The transitional epithelium may appear as a layer of column-shaped cells when relaxed, and of flatter cells when distended. Below the epithelium sits the lamina propria. The lamina propria is made up of loose connective tissue with many elastic fibers interspersed with blood vessels, veins and lymphatics. The ureter is surrounded by two muscular layers, an inner longitudinal layer of muscle, and an outer circular or spiral layer of muscle. The lower third of the ureter has a third muscular layer. Beyond these layers sits an adventitia containing blood vessels, lymphatic vessels, and veins. Development The ureters develop from the ureteric buds, which are outpouchings from the mesonephric duct. This is a duct, derived from mesoderm, found in the early embryo. Over time, the buds elongate, moving into surrounding mesodermal tissue, dilate, and divide into left and right ureters. Eventually, successive divisions from these buds form not only the ureters, but also the pelvis, major and minor calyces, and collecting ducts of the kidneys. The mesonephric duct is connected with the cloaca, which over the course of development splits into a urogenital sinus and the anorectal canal. The urinary bladder forms from the urogenital sinus. Over time, as the bladder enlarges, it absorbs the surrounding parts of the primitive ureters. Finally, the entry points of the ureters into the bladder move upwards, owing to the upward migration of the kidneys in the developing embryo. Function The ureters are a component of the urinary system. Urine, produced by the kidneys, travels along the ureters to the bladder. It does this through regular contractions called peristalsis. Clinical significance Ureteral stones A kidney stone can move from the kidney and become lodged inside the ureter, which can block the flow of urine, as well as cause a sharp cramp in the back, side, or lower abdomen. Pain often comes in waves lasting up to two hours, then subsides, called renal colic. The affected kidney could then develop hydronephrosis, should a part of the kidney become swollen due to blocked flow of urine. It is classically described that there are three sites in the ureter where a kidney stone will commonly become stuck: where the ureter meets the renal pelvis; where the iliac blood vessels cross the ureters; and where the ureters enter the urinary bladder, however a retrospective case study, which is a primary source, of where stones lodged based on medical imaging did not show many stones at the place where the iliac blood vessels cross. Most stones are compounds containing calcium such as calcium oxalate and calcium phosphate. The first recommended investigation is a CT scan of the abdomen because it can detect almost all stones. Management includes analgesia, often with nonsteroidal antiinflammatories. Small stones (< 4mm) may pass themselves; larger stones may require lithotripsy, and those with complications such as hydronephrosis or infection may require surgery to remove. Reflux Vesicoureteral reflux refers to the reflux of fluid from the bladder into the ureters. This condition can be associated with urinary tract infections, particularly in children, and is present in up to 28–36% of children to some degree. A number of forms of medical imaging are available for diagnosis of the condition, with modalities including doppler urinary tract ultrasound.Factors that affect which of these are selected depends if a child is able to receive a urinary catheter, and whether a child is toilet trained. Whether these investigations are performed at the first time a child has an illness, or later and depending on other factors (such as if the causal bacteria is E. coli) differ between US, EU and UK guidelines. Management is also variable, with differences between international guidelines on issues such as whether prophylactic antibiotics should be used, and whether surgery is recommended. One reason is most instances of vesicoureteral reflux improve by themselves. If surgery is considered, it generally involves reattaching the ureters to a different spot on the bladder, and extending the part of the ureter that it is within the wall of the bladder, with the most common surgical option being Cohen's cross-trigonal reimplantation. Anatomical and surgical abnormalities Blockage, or obstruction of the ureter can occur, as a result of narrowing within the ureter, or compression or fibrosis of structures around the ureter. Narrowing can result of ureteric stones, masses associated with cancer, and other lesions such as endometriosis tuberculosis and schistosomiasis. Things outside the ureters such as constipation and retroperitoneal fibrosis can also compress them. Some congenital abnormalities can also result in narrowing or the ureters. Congenital disorders of the ureter and urinary tract affect 10% of infants. These include partial or total duplication of the ureter (a duplex ureter), or the formation of a second irregularly placed () ureter; or where the junction with the bladder is malformed or a ureterocoele develops (usually in that location). If the ureters have been resited as a result of surgery, for example due to a kidney transplant or due to past surgery for vesicoureteric reflux, that site may also become narrowed. A narrowed ureter may lead to ureteric enlargement () and cause swelling of the kidneys (hydronephrosis). Associated symptoms may include recurrent infections, pain or blood in the urine; and when tested, kidney function might be seen to decrease. These are considered situations when surgery is needed. Medical imaging, including urinary tract ultrasound, CT or nuclear medicine imaging is conducted to investigate many causes. This may involve reinserting the ureters into a new place on the bladder (reimplantion), or widening of the ureter. A ureteric stent may be inserted to relieve an obstruction. If the cause cannot be removed, a nephrostomy may be required, which is the insertion of a tube connected to the renal pelvis which directly drains urine into a stoma bag. Cancer Cancer of the ureters is known as ureteral cancer. It is usually due to cancer of the urothelium, the cells that line the surface of the ureters. Urothelial cancer is more common after the age of 40, and more common in men than women; other risk factors include smoking and exposure to dyes such as aromatic amines and aldehydes. When cancer is present, the most common symptom is blood in the urine; it may not cause symptoms, and a physical medical examination may be otherwise normal, except in late disease. Ureteral cancer is most often due to cancer of the cells lining the ureter, called transitional cell carcinoma, although it can more rarely occur as a squamous cell carcinoma if the type of cells lining the urethra have changed due to chronic inflammation, such as due to stones or schistosomiasis. Investigations performed usually include collecting a sample of urine for an inspection for malignant cells under a microscope, called cytology, as well as medical imaging by a CT urogram or ultrasound. If a concerning lesion is seen, a flexible camera may be inserted into the ureters, called ureteroscopy, in order to view the lesion and take a biopsy, and a CT scan will be performed of other body parts (a CT scan of the chest, abdomen and pelvis) to look for additional lesions. After the cancer is staged, treatment may involve open surgery to remove the affected ureter and kidney if it is involved; or, if the lesion is small, it may be removed via ureteroscopy. Prognosis can vary markedly depending on the tumour grade, with a worse prognosis associated with an ulcerating lesion. Injury Injuries to the ureter can occur after penetrating abdominal injuries, and injuries at high speeds followed by an abrupt stop (such as a high speed car accident). The ureter can be injured during surgery to nearby structures. It is injured in 2 per 10,000 cases of vaginal hysterectomies and 13 per 10,000 cases of abdominal hysterectomies, usually near the suspensory ligament of the ovary or near the cardinal ligament, where the ureter runs close to the blood vessels of the uterus. Imaging Several forms of medical imaging are used to view the ureters and urinary tract. Ultrasound may be able to show evidence of blockage because of hydronephrosis of the kidneys and renal pelvis. CT scans, including ones where contrast media is injected intravenously to better show the ureters, and with contrast to better show lesions, and to differentiate benign from malignant lesions. Dye may also be injected directly into the ureters or renal tract; an antegrade pyelogram is when contrast is injected directly into the renal pelvis, and a retrograde pyelogram is where dye is injected into the urinary tract via a catheter, and flows backwards into the ureters. More invasive forms of imaging include ureteroscopy, which is the insertion of a flexible endoscope into the urinary tract to view the ureters. Ureteroscopy is most commonly used for medium to large-sized stones when less invasive methods of removal cannot be used. Other animals All vertebrates have two kidneys located behind the abdomen that produce urine, and have a way of excreting it, so that waste products within the urine can be removed from the body. The structure specifically called the ureter is present in amniotes, meaning mammals, birds and reptiles. These animals possess an adult kidney derived from the metanephros. The duct that connects the kidney to excrete urine in these animals is the ureter. In placental mammals, it connects to the urinary bladder, whence urine leaves via the urethra. In monotremes, urine flows from the ureters into the cloaca. The ureters are ventral to the vasa deferentia in male placental mammals, but dorsal to the vasa deferentia in marsupials. In female marsupials, the ureters pass between the median and lateral vaginae. History The word "ureter" comes from the Ancient Greek noun , , meaning "urine", and the first use of the word is seen during the era of Hippocrates to refer to the urethra. The anatomical structure of the ureter was noted by 40 AD. However, the terms "ureter" and "urethra" were variably used to refer to each other thereafter for more than a millennium. It was only in the 1550s that anatomists such as Bartolomeo Eustachi and Jacques Dubois began to use the terms to specifically and consistently refer to what are in modern English called the ureter and the urethra. Following this, in the 19th and 20th centuries, multiple terms relating to the structures such as ureteritis and ureterography, were coined. Kidney stones have been identified and recorded about as long as written historical records exist. The urinary tract including the ureters, as well as their function to drain urine from the kidneys, has been described by Galen in the second century AD. The first to examine the ureter through an internal approach, called ureteroscopy, rather than surgery was Hampton Young in 1929. This was improved on by VF Marshall who is the first published use of a flexible endoscope based on fiber optics, which occurred in 1964. The insertion of a drainage tube into the renal pelvis, bypassing the ureters and urinary tract, called nephrostomy, was first described in 1941. Such an approach differed greatly from the open surgical approaches within the urinary system employed during the preceding two millennia. The first radiological imaging of the ureters was by X-rays, although this was made more difficult by the thick abdomen, which the low power of the original X-rays could not penetrate enough to produce clear images. More useful images were able to be produced when Edwin Hurry Fenwick in 1908 pioneered the use of tubes covered in material visible to X-rays inserted into the ureters, and in the early 20th century when contrasts were injected externally into the urinary tract (retrograde pyelograms). Unfortunately, much of the earlier retrograde pyelograms were complicated by significant damage to the kidneys as a result of contrast based on silver or sodium iodide. Hryntshalk in 1929 pioneered the development of the intravenous urogram, in which contrast is injected into a vein and highlights the kidney and, when excreted, the urinary tract. Things improved with the development by Moses Swick and Leopold Lichtwitz in the late 1920s of relatively nontoxic contrast media, with controversy surrounding publication as to who was the primary discoverer. Side-effects associated with imaging improved even more when Tosten Almen published a ground-breaking thesis in 1969 based on the less toxic low-osmolar contrast media, developed based on swimming experiences in lakes with different salinity.	Biology and health sciences	Urinary system	Biology
188378	https://en.wikipedia.org/wiki/Phosphoric%20acid	Phosphoric acid	Phosphoric acid (orthophosphoric acid, monophosphoric acid or phosphoric(V) acid) is a colorless, odorless phosphorus-containing solid, and inorganic compound with the chemical formula . It is commonly encountered as an 85% aqueous solution, which is a colourless, odourless, and non-volatile syrupy liquid. It is a major industrial chemical, being a component of many fertilizers. The compound is an acid. Removal of all three ions gives the phosphate ion . Removal of one or two protons gives dihydrogen phosphate ion , and the hydrogen phosphate ion , respectively. Phosphoric acid forms esters, called organophosphates. The name "orthophosphoric acid" can be used to distinguish this specific acid from other "phosphoric acids", such as pyrophosphoric acid. Nevertheless, the term "phosphoric acid" often means this specific compound; and that is the current IUPAC nomenclature. Production Phosphoric acid is produced industrially by one of two routes, wet processes and dry. Wet process In the wet process, a phosphate-containing mineral such as calcium hydroxyapatite and fluorapatite are treated with sulfuric acid. Calcium sulfate (gypsum, ) is a by-product, which is removed as phosphogypsum. The hydrogen fluoride (HF) gas is streamed into a wet (water) scrubber producing hydrofluoric acid. In both cases the phosphoric acid solution usually contains 23–33% (32–46% ). It may be concentrated to produce commercial- or merchant-grade phosphoric acid, which contains about 54–62% (75–85% ). Further removal of water yields superphosphoric acid with a concentration above 70% (corresponding to nearly 100% ). The phosphoric acid from both processes may be further purified by removing compounds of arsenic and other potentially toxic impurities. Dry process To produce food-grade phosphoric acid, phosphate ore is first reduced with coke in an electric arc furnace, to give elemental phosphorus. This process is also known as the thermal process or the electric furnace process. Silica is also added, resulting in the production of calcium silicate slag. Elemental phosphorus is distilled out of the furnace and burned with air to produce high-purity phosphorus pentoxide, which is dissolved in water to make phosphoric acid. The thermal process produces phosphoric acid with a very high concentration of (about 85%) and a low level of impurities. However, this process is more expensive and energy-intensive than the wet process, which produces phosphoric acid with a lower concentration of (about 26-52%) and a higher level of impurities. The wet process is the most common method of producing phosphoric acid for fertilizer use. Even in China, where the thermal process is still used quite widely due to relatively cheap coal as opposed to the sulfuric acid, over 7/8 of phosphoric acid is produced with wet process. Purification Phosphoric acids produced from phosphate rock or thermal processes often requires purification. A common purification methods is liquid-liquid extraction, which involves the separation of phosphoric acids from water and other impurities using organic solvents, such as tributyl phosphate (TBP), methyl isobutyl ketone (MIBK), or n-octanol. Nanofiltration involves the use of a premodified nanofiltration membrane, which is functionalized by a deposit of a high molecular weight polycationic polymer of polyethyleneimines. Nanofiltration has been shown to significantly reduce the concentrations of various impurities, including cadmium, aluminum, iron, and rare earth elements. The laboratory and industrial pilot scale results showed that this process allows the production of food-grade phosphoric acid. Fractional crystallization can achieve highest purities typically used for semiconductor applications. Usually a static crystallizer is used. A static crystallizer uses vertical plates, which are suspended in the molten feed and which are alternatingly cooled and heated by a heat transfer medium. The process begins with the slow cooling of the heat transfer medium below the freezing point of the stagnant melt. This cooling causes a layer of crystals to grow on the plates. Impurities are rejected from the growing crystals and are concentrated in the remaining melt. After the desired fraction has been crystallized, the remaining melt is drained from the crystallizer. The purer crystalline layer remains adhered to the plates. In a subsequent step, the plates are heated again to liquify the crystals and the purified phosphoric acid drained into the product vessel. The crystallizer is filled with feed again and the next cooling cycle is started. Properties Acidic properties In aqueous solution phosphoric acid behaves as a triprotic acid. , pKa1 = 2.14 , pKa2 = 7.20 , pKa3 = 12.37 The difference between successive pKa values is sufficiently large so that salts of either monohydrogen phosphate, or dihydrogen phosphate, , can be prepared from a solution of phosphoric acid by adjusting the pH to be mid-way between the respective pKa values. Aqueous solutions Aqueous solutions up to 62.5% are eutectic, exhibiting freezing-point depression as low as -85°C. When the concentration of acid rises above 62.5% the freezing-point increases, reaching 21°C by 85% (w/w; the monohydrate). Beyond this the phase diagram becomes complicated, with significant local maxima and minima. For this reason phosphoric acid is rarely sold above 85%, as beyond this adding or removing small amounts moisture risks the entire mass freezing solid, which would be a major problem on a large scale. A local maximum at 91.6% which corresponds to the hemihydrate 2H3PO4•H2O, freezing at 29.32°C. There is a second smaller eutectic depression at a concentration of 94.75% with a freezing point of 23.5°C. At higher concentrations the freezing point rapidly increases. Concentrated phosphoric acid tends to supercool before crystallization occurs, and may be relatively resistant to crystallisation even when stored below the freezing point. Self condensation Phosphoric acid is commercially available as aqueous solutions of various concentrations, not usually exceeding 85%. If concentrated further it undergoes slow self-condensation, forming an equilibrium with pyrophosphoric acid: Even at 90% concentration the amount of pyrophosphoric acid present is negligible, but beyond 95% it starts to increase, reaching 15% at what would have otherwise been 100% orthophosphoric acid. As the concentration is increased higher acids are formed, culminating in the formation of polyphosphoric acids. It is not possible to fully dehydrate phosphoric acid to phosphorus pentoxide, instead the polyphosphoric acid becomes increasingly polymeric and viscous. Due to the self-condensation, pure orthophosphoric acid can only be obtained by a careful fractional freezing/melting process. Uses The dominant use of phosphoric acid is for fertilizers, consuming approximately 90% of production. Food-grade phosphoric acid (additive E338) is used to acidify foods and beverages such as various colas and jams, providing a tangy or sour taste. The phosphoric acid also serves as a preservative. Soft drinks containing phosphoric acid, which would include Coca-Cola, are sometimes called phosphate sodas or phosphates. Phosphoric acid in soft drinks has the potential to cause dental erosion. Phosphoric acid also has the potential to contribute to the formation of kidney stones, especially in those who have had kidney stones previously. Specific applications of phosphoric acid include: in anti-rust treatment by phosphate conversion coating or passivation to prevent iron oxidation by means of the Parkerization process as an external standard for phosphorus-31 nuclear magnetic resonance in phosphoric acid fuel cells in activated carbon production in compound semiconductor processing, to etch Indium gallium arsenide selectively with respect to indium phosphide in microfabrication to etch silicon nitride selectively with respect to silicon dioxide in microfabrication to etch aluminium as a pH adjuster in cosmetics and skin-care products as a sanitizing agent in the dairy, food, and brewing industries Phosphoric acid may also be used for chemical polishing (etching) of metals like aluminium or for passivation of steel products in a process called phosphatization. Safety Phosphoric acid is not a strong acid. However, at moderate concentrations phosphoric acid solutions are irritating to the skin. Contact with concentrated solutions can cause severe skin burns and permanent eye damage. A link has been shown between long-term regular cola intake and osteoporosis in later middle age in women (but not men).	Physical sciences	Inorganic compounds	null
188386	https://en.wikipedia.org/wiki/Lighting	Lighting	Lighting or illumination is the deliberate use of light to achieve practical or aesthetic effects. Lighting includes the use of both artificial light sources like lamps and light fixtures, as well as natural illumination by capturing daylight. Daylighting (using windows, skylights, or light shelves) is sometimes used as the main source of light during daytime in buildings. This can save energy in place of using artificial lighting, which represents a major component of energy consumption in buildings. Proper lighting can enhance task performance, improve the appearance of an area, or have positive psychological effects on occupants. Indoor lighting is usually accomplished using light fixtures, and is a key part of interior design. Lighting can also be an intrinsic component of landscape projects. History With the discovery of fire, the earliest form of artificial lighting used to illuminate an area were campfires or torches. As early as 400,000 years ago, fire was kindled in the caves of Peking Man. Prehistoric people used primitive oil lamps to illuminate surroundings. These lamps were made from naturally occurring materials such as rocks, shells, horns and stones, were filled with grease, and had a fiber wick. Lamps typically used animal or vegetable fats as fuel. Hundreds of these lamps (hollow worked stones) have been found in the Lascaux caves in modern-day France, dating to about 15,000 years ago. Oily animals (birds and fish) were also used as lamps after being threaded with a wick. Fireflies have been used as lighting sources. Candles and glass and pottery lamps were also invented. Chandeliers were an early form of "light fixture". A major reduction in the cost of lighting occurred with the discovery of whale oil. The use of whale oil declined after Abraham Gesner, a Canadian geologist, first refined kerosene in the 1840s, allowing brighter light to be produced at substantially lower cost. In the 1850s, the price of whale oil dramatically increased (more than doubling from 1848 to 1856) due to shortages of available whales, hastening whale oil's decline. By 1860, there were 33 kerosene plants in the United States, and Americans spent more on gas and kerosene than on whale oil. The final death knell for whale oil was in 1859, when crude oil was discovered and the petroleum industry arose. Gas lighting was economical enough to power street lights in major cities starting in the early 1800s, and was also used in some commercial buildings and in the homes of wealthy people. The gas mantle boosted the luminosity of utility lighting and of kerosene lanterns. The next major drop in price came about in the 1880s with the introduction of electric lighting in the form of arc lights for large space and street lighting, followed by incandescent light bulb-based utilities for indoor and outdoor lighting. Over time, electric lighting became ubiquitous in developed countries. Segmented sleep patterns disappeared, improved nighttime lighting made more activities possible at night, and more street lights reduced urban crime. Fixtures Lighting fixtures come in a wide variety of styles for various functions. The most important functions are as a holder for the light source, to provide directed light and to avoid visual glare. Some are very plain and functional, while some are pieces of art in themselves. Nearly any material can be used, so long as it can tolerate the excess heat and is in keeping with safety codes. An important property of light fixtures is the luminous efficacy or wall-plug efficiency, meaning the amount of usable light emanating from the fixture per used energy, usually measured in lumen per watt. A fixture using replaceable light sources can also have its efficiency quoted as the percentage of light passed from the "bulb" to the surroundings. The more transparent the lighting fixtures are, the higher efficacy. Shading the light will normally decrease efficacy but increase the directionality and the visual comfort probability. Color temperature for white light sources also affects their use for certain applications. The color temperature of a white light source is the temperature in kelvins of a theoretical black body emitter that most closely matches the spectral characteristics (spectral power distribution) of the lamp. An incandescent bulb has a color temperature around 2800 to 3000 kelvins; daylight is around 6400 kelvins. Lower color temperature lamps have relatively more energy in the yellow and red part of the visible spectrum, while high color temperatures correspond to lamps with more of a blue-white appearance. For critical inspection or color matching tasks, or for retail displays of food and clothing, the color temperature of the lamps will be selected for the best overall lighting effect. Types Lighting is classified by intended use as general, accent, or task lighting, depending largely on the distribution of the light produced by the fixture. Task lighting is mainly functional and is usually the most concentrated, for purposes such as reading or inspection of materials. For example, reading poor-quality reproductions may require task lighting levels up to 1500 lux (140 footcandles), and some inspection tasks or surgical procedures require even higher levels. Accent lighting is mainly decorative, intended to highlight pictures, plants, or other elements of interior design or landscaping. General lighting (sometimes referred to as ambient light) fills in between the two and is intended for general illumination of an area. Indoors, this would be a basic lamp on a table or floor, or a fixture on the ceiling. Outdoors, general lighting for a parking lot may be as low as 10-20 lux (1-2 footcandles) since pedestrians and motorists already used to the dark will need little light for crossing the area. Methods Downlighting is most common, with fixtures on or recessed in the ceiling casting light downward. This tends to be the most used method, used in both offices and homes. Although it is easy to design, it has dramatic problems with glare and excess energy consumption due to large number of fittings. The introduction of LED lighting has greatly improved this by approx. 90% when compared to a halogen downlight or spotlight. LED lamps or bulbs are now available to retro fit in place of high energy consumption lamps. Uplighting is less common, often used to bounce indirect light off the ceiling and back down. It is commonly used in lighting applications that require minimal glare and uniform general illuminance levels. Uplighting (indirect) uses a diffuse surface to reflect light in a space and can minimize disabling glare on computer displays and other dark glossy surfaces. It gives a more uniform presentation of the light output in operation. However indirect lighting is completely reliant upon the reflectance value of the surface. While indirect lighting can create a diffused and shadow free light effect it can be regarded as an uneconomical lighting principle. Front lighting is also quite common, but tends to make the subject look flat as its casts almost no visible shadows. Lighting from the side is the less common, as it tends to produce glare near eye level. Backlighting either around or through an object is mainly for accent. Backlighting is used to illuminate a background or backdrop. This adds depth to an image or scene. Others use it to achieve a more dramatic effect. Forms of lighting Indoor lighting Forms of lighting include alcove lighting, which like most other uplighting is indirect. This is often done with fluorescent lighting (first available at the 1939 World's Fair) or rope light, occasionally with neon lighting, and recently with LED strip lighting. It is a form of backlighting. Soffit or close to wall lighting can be general or a decorative wall-wash, sometimes used to bring out texture (like stucco or plaster) on a wall, though this may also show its defects as well. The effect depends heavily on the exact type of lighting source used. Recessed lighting (often called "pot lights" in Canada, "can lights" or 'high hats" in the US) is popular, with fixtures mounted into the ceiling structure so as to appear flush with it. These downlights can use narrow beam spotlights, or wider-angle floodlights, both of which are bulbs having their own reflectors. There are also downlights with internal reflectors designed to accept common 'A' lamps (light bulbs) which are generally less costly than reflector lamps. Downlights can be incandescent, fluorescent, HID (high intensity discharge) or LED. Track lighting, invented by Lightolier, was popular at one period of time because it was much easier to install than recessed lighting, and individual fixtures are decorative and can be easily aimed at a wall. It has regained some popularity recently in low-voltage tracks, which often look nothing like their predecessors because they do not have the safety issues that line-voltage systems have, and are therefore less bulky and more ornamental in themselves. A master transformer feeds all of the fixtures on the track or rod with 12 or 24 volts, instead of each light fixture having its own line-to-low voltage transformer. There are traditional spots and floods, as well as other small hanging fixtures. A modified version of this is cable lighting, where lights are hung from or clipped to bare metal cables under tension. A sconce is a wall-mounted fixture, particularly one that shines up and sometimes down as well. A torchère is an uplight intended for ambient lighting. It is typically a floor lamp but may be wall-mounted like a sconce. Further interior light fixtures include chandeliers, pendant lights, ceiling fans with lights, close-to-ceiling or flush lights, and various types of lamps The portable or table lamp is probably the most common fixture, found in many homes and offices. The standard lamp and shade that sits on a table is general lighting, while the desk lamp is considered task lighting. Magnifier lamps are also task lighting. The illuminated ceiling was once popular in the 1960s and 1970s but fell out of favor after the 1980s. This uses diffuser panels hung like a suspended ceiling below fluorescent lights, and is considered general lighting. Other forms include neon, which is not usually intended to illuminate anything else, but to actually be an artwork in itself. This would probably fall under accent lighting, though in a dark nightclub it could be considered general lighting. In a movie theater, steps in the aisles are usually marked with a row of small lights for convenience and safety, when the film has started and the other lights are off. Traditionally made up of small low wattage, low-voltage lamps in a track or translucent tube, these are rapidly being replaced with LED based versions. Outdoor lighting Street Lights are used to light roadways and walkways at night. Some manufacturers are designing LED and photovoltaic luminaires to provide an energy-efficient alternative to traditional street light fixtures. Floodlights can be used to illuminate work zones or outdoor playing fields during nighttime hours. The most common type of floodlights are metal halide and high pressure sodium lights. Beacon lights are positioned at the intersection of two roads to aid in navigation. Sometimes security lighting can be used along roadways in urban areas, or behind homes or commercial facilities. These are extremely bright lights used to deter crime. Security lights may include floodlights and be activated with PIR switches that detect moving heat sources in darkness. Entry lights can be used outside to illuminate and signal the entrance to a property. These lights are installed for safety, security, and for decoration. Underwater accent lighting is also used for koi ponds, fountains, swimming pools and the like. Neon signs are most often used to attract attention rather than to illuminate. Vehicle use Vehicles typically include headlamps and tail lights. Headlamps are white or selective yellow lights placed in the front of the vehicle, designed to illuminate the upcoming road and to make the vehicle more visible. Many manufactures are turning to LED headlights as an energy-efficient alternative to traditional headlamps. Tail and brake lights are red and emit light to the rear so as to reveal the vehicle's direction of travel to following drivers. White rear-facing reversing lamps indicate that the vehicle's transmission has been placed in the reverse gear, warning anyone behind the vehicle that it is moving backwards, or about to do so. Flashing turn signals on the front, side, and rear of the vehicle indicate an intended change of position or direction. In the late 1950s, some automakers began to use electroluminescent technology to backlight their cars' speedometers and other gauges or to draw attention to logos or other decorative elements. Lamps Commonly called 'light bulbs', lamps are the removable and replaceable part of a light fixture, which converts electrical energy into electromagnetic radiation. While lamps have traditionally been rated and marketed primarily in terms of their power consumption, expressed in watts, proliferation of lighting technology beyond the incandescent light bulb has eliminated the correspondence of wattage to the amount of light produced. For example, a 60 W incandescent light bulb produces about the same amount of light as a 13 W compact fluorescent lamp. Each of these technologies has a different efficacy in converting electrical energy to visible light. Visible light output is typically measured in lumens. This unit only quantifies the visible radiation, and excludes invisible infrared and ultraviolet light. A wax candle produces on the close order of 13 lumens, a 60 watt incandescent lamp makes around 700 lumens, and a 15-watt compact fluorescent lamp produces about 800 lumens, but actual output varies by specific design. Rating and marketing emphasis is shifting away from wattage and towards lumen output, to give the purchaser a directly applicable basis upon which to select a lamp. Lamp types include: Ballast: A ballast is an auxiliary piece of equipment designed to start and properly control the flow of power to discharge light sources such as fluorescent and high intensity discharge (HID) lamps. Some lamps require the ballast to have thermal protection. fluorescent light: A tube coated with phosphor containing low pressure mercury vapor that produces white light. Halogen: Incandescent lamps containing halogen gases such as iodine or bromine, increasing the efficacy of the lamp versus a plain incandescent lamp. Neon: A low pressure gas contained within a glass tube; the color emitted depends on the gas. Light-emitting diodes: Light-emitting diodes (LED) are solid state devices that emit light by dint of the movement of electrons in a semiconductor material. Compact fluorescent lamps: CFLs are designed to replace incandescent lamps in existing and new installations. Design and architecture Architectural lighting design Lighting design as it applies to the built environment is known as 'architectural lighting design'. Lighting of structures considers aesthetic elements as well as practical considerations of quantity of light required, occupants of the structure, energy efficiency, and cost. Artificial lighting takes into account the amount of daylight received in a space by using daylight factor calculations. For simple installations, hand calculations based on tabular data are used to provide an acceptable lighting design. More critical or complex designs now routinely use computer software such as Radiance for mathematical modeling, which can allow an architect to quickly and accurately evaluate the benefit of a proposed design. In some instances, the materials used on walls and furniture play a key role in the lighting effect. For example, dark paint tends to absorb light, making the room appear smaller and more dim than it is, whereas light paint does the opposite. Other reflective surfaces also have an effect on lighting design. On stage and set Lighting illuminates the performers and artists in a live theatre, dance, or musical performance, and is selected and arranged to create dramatic effects. Stage lighting uses general illumination technology in devices configured for easy adjustment of their output characteristics. The setup of stage lighting is tailored for each scene of each production. Dimmers, colored filters, reflectors, lenses, motorized or manually aimed lamps, and different kinds of flood and spot lights are among the tools used by a stage lighting designer to produce the desired effects. A set of lighting cues are prepared so that the lighting operator can control the lights in step with the performance; complex theatre lighting systems use computer control of lighting instruments. Motion picture and television production use many of the same tools and methods of stage lighting. Especially in the early days of these industries, very high light levels were required and heat produced by lighting equipment presented substantial challenges. Modern cameras require less light, and modern light sources emit less heat. Measurement Measurement of light or photometry is generally concerned with the amount of useful light falling on a surface and the amount of light emerging from a lamp or other source, along with the colors that can be rendered by this light. The human eye responds differently to light from different parts of the visible spectrum, therefore photometric measurements must take the luminosity function into account when measuring the amount of useful light. The basic SI unit of measurement is the candela (cd), which describes the luminous intensity, all other photometric units are derived from the candela. Luminance for instance is a measure of the density of luminous intensity in a given direction. It describes the amount of light that passes through or is emitted from a particular area, and falls within a given solid angle. The SI unit for luminance is candela per square metre (cd/m2). The CGS unit of luminance is the stilb, which is equal to one candela per square centimetre or 10 kcd/m2. The amount of useful light emitted from a source or the luminous flux is measured in lumen (lm). The SI unit of illuminance and luminous emittance, being the luminous power per area, is measured in Lux. It is used in photometry as a measure of the intensity, as perceived by the human eye, of light that hits or passes through a surface. It is analogous to the radiometric unit watts per square metre, but with the power at each wavelength weighted according to the luminosity function, a standardized model of human visual brightness perception. In English, "lux" is used in both singular and plural. Visual comfort often entails the measurement of subjective evaluations. Several measurement methods have been developed to control glare resulting from indoor lighting design. The Unified Glare Rating (UGR), the Visual Comfort Probability, and the Daylight Glare Index are some of the most well-known methods of measurement. In addition to these new methods, four main factors influence the degree of discomfort glare; the luminance of the glare source, the solid angle of the glare source, the background luminance, and the position of the glare source in the field of view must all be taken into account. Color properties To define light source color properties, the lighting industry predominantly relies on two metrics, correlated color temperature (CCT), commonly used as an indication of the apparent "warmth" or "coolness" of the light emitted by a source, and color rendering index (CRI), an indication of the light source's ability to make objects appear natural. However, these two metrics, developed in the last century, are facing increased challenges and criticisms as new types of light sources, particularly light-emitting diodes (LEDs), become more prevalent in the market. For example, in order to meet the expectations for good color rendering in retail applications, research suggests using the well-established CRI along with another metric called gamut area index (GAI). GAI represents the relative separation of object colors illuminated by a light source; the greater the GAI, the greater the apparent saturation or vividness of the object colors. As a result, light sources which balance both CRI and GAI are generally preferred over ones that have only high CRI or only high GAI. Light exposure Typical measurements of light have used a Dosimeter. Dosimeters measure an individual's or an object's exposure to something in the environment, such as light dosimeters and ultraviolet dosimeters. In order to specifically measure the amount of light entering the eye, personal circadian light meter called the Daysimeter has been developed. This is the first device created to accurately measure and characterize light (intensity, spectrum, timing, and duration) entering the eye that affects the human body's clock. The small, head-mounted device measures an individual's daily rest and activity patterns, as well as exposure to short-wavelength light that stimulates the circadian system. The device measures activity and light together at regular time intervals and electronically stores and logs its operating temperature. The Daysimeter can gather data for up to 30 days for analysis. Energy consumption Several strategies are available to minimize energy requirements for lighting a building: Specification of illumination requirements for each given use area Analysis of lighting quality to ensure that adverse components of lighting (for example, glare or incorrect color spectrum) are not biasing the design Integration of space planning and interior architecture (including choice of interior surfaces and room geometries) to lighting design Design of time of day use that does not expend unnecessary energy Selection of fixtures and lamps that reflect best available technology for energy conservation Training of building occupants to use lighting equipment in most efficient manner Maintenance of lighting systems to minimize energy wastage Use of natural light Some big box stores were being built from 2006 on with numerous plastic bubble skylights, in many cases completely obviating the need for interior artificial lighting for many hours of the day. In countries where indoor lighting of simple dwellings is a significant cost, "Moser lamps", plastic water-filled transparent drink bottles fitted through the roof, provide the equivalent of a 40- to 60-watt incandescent bulb each during daylight. Load shedding can help reduce the power requested by individuals to the main power supply. Load shedding can be done on an individual level, at a building level, or even at a regional level. Specification of illumination requirements is the basic concept of deciding how much illumination is required for a given task. Clearly, much less light is required to illuminate a hallway compared to that needed for a word processing work station. Generally speaking, the energy expended is proportional to the design illumination level. For example, a lighting level of 400 lux might be chosen for a work environment involving meeting rooms and conferences, whereas a level of 80 lux could be selected for building hallways. If the hallway standard simply emulates the conference room needs, then much more energy will be consumed than is needed. Lighting control systems Lighting control systems reduce energy usage and cost by helping to provide light only when and where it is needed. Lighting control systems typically incorporate the use of time schedules, occupancy control, and photocell control (i.e. daylight harvesting). Some systems also support demand response and will automatically dim or turn off lights to take advantage of utility incentives. Lighting control systems are sometimes incorporated into larger building automation systems. Many newer control systems are using wireless mesh open standards (such as Zigbee), which provides benefits including easier installation (no need to run control wires) and interoperability with other standards-based building control systems (e.g. security). In response to daylighting technology, daylight harvesting systems have been developed to further reduce energy consumption. These technologies are helpful, but they do have their downfalls. Many times, rapid and frequent switching of the lights on and off can occur, particularly during unstable weather conditions or when daylight levels are changing around the switching illuminance. Not only does this disturb occupants, it can also reduce lamp life. A variation of this technology is the 'differential switching or dead-band' photoelectric control which has multiple illuminances it switches from so as not to disturb occupants as much. Occupancy sensors to allow operation for whenever someone is within the area being scanned can control lighting. When motion can no longer be detected, the lights shut off. Passive infrared sensors react to changes in heat, such as the pattern created by a moving person. The control must have an unobstructed view of the building area being scanned. Doors, partitions, stairways, etc. will block motion detection and reduce its effectiveness. The best applications for passive infrared occupancy sensors are open spaces with a clear view of the area being scanned. Ultrasonic sensors transmit sound above the range of human hearing and monitor the time it takes for the sound waves to return. A break in the pattern caused by any motion in the area triggers the control. Ultrasonic sensors can see around obstructions and are best for areas with cabinets and shelving, restrooms, and open areas requiring 360-degree coverage. Some occupancy sensors utilize both passive infrared and ultrasonic technology, but are usually more expensive. They can be used to control one lamp, one fixture or many fixtures. Daylighting Daylighting is the oldest method of interior lighting. Daylighting is simply designing a space to use as much natural light as possible. This decreases energy consumption and costs, and requires less heating and cooling from the building. Daylighting has also been proven to have positive effects on patients in hospitals as well as work and school performance. Due to a lack of information that indicate the likely energy savings, daylighting schemes are not yet popular among most buildings. Unlike electric lighting, the distribution of daylight varies considerably throughout the entire year inside a building. Solid-state lighting In recent years light-emitting diodes (LEDs) are becoming increasingly efficient leading to an extraordinary increase in the use of solid state lighting. In many situations, controlling the light emission of LEDs may be done most effectively by using the principles of nonimaging optics. Health effects It is valuable to provide the correct light intensity and color spectrum for each task or environment. Otherwise, energy not only could be wasted but over-illumination can lead to adverse health and psychological effects. Beyond the energy factors being considered, it is important not to over-design illumination, lest adverse health effects such as headache frequency, stress, and increased blood pressure be induced by the higher lighting levels. In addition, glare or excess light can decrease worker efficiency. Analysis of lighting quality particularly emphasizes use of natural lighting, but also considers spectral content if artificial light is to be used. Not only will greater reliance on natural light reduce energy consumption, but will favorably impact human health and performance. New studies have shown that the performance of students is influenced by the time and duration of daylight in their regular schedules. Designing school facilities to incorporate the right types of light at the right time of day for the right duration may improve student performance and well-being. Similarly, designing lighting systems that maximize the right amount of light at the appropriate time of day for the elderly may help relieve symptoms of Alzheimer's disease. The human circadian system is entrained to a 24-hour light-dark pattern that mimics the earth's natural light/dark pattern. When those patterns are disrupted, they disrupt the natural circadian cycle. Circadian disruption may lead to numerous health problems including breast cancer, seasonal affective disorder, delayed sleep phase syndrome, and other ailments. A study conducted in 1972 and 1981, documented by Robert Ulrich, surveyed 23 surgical patients assigned to rooms looking out on a natural scene. The study concluded that patients assigned to rooms with windows allowing much natural light had shorter postoperative hospital stays, received fewer negative evaluative comments in nurses' notes, and took fewer potent analgesics than 23 matched patients in similar rooms with windows facing a brick wall. This study suggests that due to the nature of the scenery and daylight exposure was indeed healthier for patients as opposed to those exposed to little light from the brick wall. In addition to increased work performance, proper usage of windows and daylighting crosses the boundaries between pure aesthetics and overall health. Alison Jing Xu, assistant professor of management at the University of Toronto Scarborough and Aparna Labroo of Northwestern University conducted a series of studies analyzing the correlation between lighting and human emotion. The researchers asked participants to rate a number of things such as: the spiciness of chicken-wing sauce, the aggressiveness of a fictional character, how attractive someone was, their feelings about specific words, and the taste of two juices–all under different lighting conditions. In their study, they found that both positive and negative human emotions are felt more intensely in bright light. Professor Xu stated, "we found that on sunny days depression-prone people actually become more depressed." They also found that dim light makes people make more rational decisions and settle negotiations easier. In the dark, emotions are slightly suppressed. However, emotions are intensified in the bright light. Environmental issues Compact fluorescent lamps Compact fluorescent lamps (CFLs) use less power than an incandescent lamp to supply the same amount of light, however they contain mercury which is a disposal hazard. Due to the ability to reduce electricity consumption, many organizations encouraged the adoption of CFLs. Some electric utilities and local governments subsidized CFLs or provided them free to customers as a means of reducing electricity demand. For a given light output, CFLs use between one fifth and one quarter the power of an equivalent incandescent lamp. Unlike incandescent lamps CFLs need a little time to warm up and reach full brightness. Not all CFLs are suitable for dimming. CFL's have largely been replaced with LED technologies. LED lamps LED lamps provide significant energy savings over incandescent and fluorescent lamps. According to the Energy Saving Trust, LED lamps use only 10% power compared to a standard incandescent bulb, where compact fluorescent lamps use 20% and energy saving halogen lamps 70%. The lifetime is also much longer — up to 50,000 hours. The downside when they were first popularized was the initial cost. By 2018, production costs dropped, performance increased, and energy consumption was reduced. While the initially cost of LEDs is still higher than incandescent lamps, the savings are so dramatic that there are very few instances that LEDs are not the most economical choice. Scattered light from outdoor illumination may have effects on the environment and human health. For instance, one study conducted by the American Medical Association warned on the use of high blue content white LEDs in street lighting, due to their higher impact on human health and environment, compared to low blue content light sources (e.g. High Pressure Sodium, phosphor-coated or PC amber LEDs, and low CCT LEDs). Light pollution Light pollution is a growing problem in reaction to excess light being given off by numerous signs, houses, and buildings. Polluting light is often wasted light involving unnecessary energy costs and carbon dioxide emissions. Light pollution is described as artificial light that is excessive or intrudes where it is not wanted. Well-designed lighting sends light only where it is needed without scattering it elsewhere. Poorly designed lighting can also compromise safety. For example, glare creates safety issues around buildings by causing very sharp shadows, temporarily blinding passersby making them vulnerable to would-be assailants. The negative ecological effects of artificial light have been increasingly well documented. The World Health Organization in 2007 issued a report that noted the effects of bright light on flora and fauna, sea turtle hatchlings, frogs during mating season and the migratory patterns of birds. The American Medical Association in 2012 issued a warning that extended exposure to light at night increases the risk of some cancers. Two studies in Israel from 2008 have yielded some additional findings about a possible correlation between artificial light at night and certain cancers. Effects on animals Artificial light at night refers to any light source other than a natural light source. Sources of artificial light include LEDS and fluorescents. This particular light source has effect on the reproduction, immune function, metabolism, thermoregulation and body temperature of organisms that need light for their daily activity. Firstly, most organisms metabolism largely depends on light. In some instances the presence of intense light starts up or increases enzyme activity inside the body of an animal. For diurnal organisms, high rate of metabolism takes place during the day and reduces or comes to a stop during the night thus, artificial light at night has a negative impact of the metabolism of diurnal organisms. Moreover, the body temperature of diurnal animals fall during the night but the presence of artificial light at night, then causes an increase in body temperature which affects the melatonin levels of the animal. Furthermore, for organisms such as birds, their sex organs are activated in relation to light intensity in certain periods during the summer at day time to aid reproduction. These sex organs are deactivated during the night but the presence of artificial light during the night sometimes disrupts their reproduction process. Professional organizations International The International Commission on Illumination (CIE) is an international authority and standard defining organization on color and lighting. Publishing widely used standard metrics such as various CIE color spaces and the color rendering index. The Illuminating Engineering Society (IES), in conjunction with organizations like ANSI and ASHRAE, publishes guidelines, standards, and handbooks that allow categorization of the illumination needs of different built environments. Manufacturers of lighting equipment publish photometric data for their products, which defines the distribution of light released by a specific luminaire. This data is typically expressed in standardized form defined by the IESNA. The International Association of Lighting Designers (IALD) is an organization which focuses on the advancement of lighting design education and the recognition of independent professional lighting designers. Those fully independent designers who meet the requirements for professional membership in the association typically append the abbreviation IALD to their name. The Professional Lighting Designers Association (PLDA), formerly known as ELDA is an organisation focusing on the promotion of the profession of Architectural Lighting Design. They publish a monthly newsletter and organise different events throughout the world. The National Council on Qualifications for the Lighting Professions (NCQLP) offers the Lighting Certification Examination which tests rudimentary lighting design principles. Individuals who pass this exam become "Lighting Certified" and may append the abbreviation LC to their name. This certification process is one of three national (U.S.) examinations (the others are CLEP and CLMC) in the lighting industry and is open not only to designers, but to lighting equipment manufacturers, electric utility employees, etc. The Professional Lighting And Sound Association (PLASA) is a UK-based trade organisation representing the 500+ individual and corporate members drawn from the technical services sector. Its members include manufacturers and distributors of stage and entertainment lighting, sound, rigging and similar products and services, and affiliated professionals in the area. They lobby for and represent the interests of the industry at various levels, interacting with government and regulating bodies and presenting the case for the entertainment industry. Example subjects of this representation include the ongoing review of radio frequencies (which may or may not affect the radio bands in which wireless microphones and other devices use) and engaging with the issues surrounding the introduction of the RoHS (Restriction of Hazardous Substances Directive) regulations. National Association de Concepteurs Eclairage (ACE) in France American Lighting Association (ALA) in the United States Associazione Professionisti dell'Illuminazione (APIL) in Italy Hellenic Illumination Committee (HIC) in Greece Indian Society of Lighting Engineers (ISLE) Institution of Lighting Engineers (ILE) in the United Kingdom Schweizerische Licht Gesellschaft (SLG) in Switzerland Society of Light and Lighting (SLL), part of the Chartered Institution of Building Services Engineers in the United Kingdom. United Scenic Artists Local 829 (USA829), membership for lighting designers as a category, with scenic designers, projection designers, costume designers, and sound designers, in the United States	Technology	Architectural elements	null
188401	https://en.wikipedia.org/wiki/Axiomatic%20system	Axiomatic system	In mathematics and logic, an axiomatic system is any set of primitive notions and axioms to logically derive theorems. A theory is a consistent, relatively-self-contained body of knowledge which usually contains an axiomatic system and all its derived theorems. An axiomatic system that is completely described is a special kind of formal system. A formal theory is an axiomatic system (usually formulated within model theory) that describes a set of sentences that is closed under logical implication. A formal proof is a complete rendition of a mathematical proof within a formal system. Properties An axiomatic system is said to be consistent if it lacks contradiction. That is, it is impossible to derive both a statement and its negation from the system's axioms. Consistency is a key requirement for most axiomatic systems, as the presence of contradiction would allow any statement to be proven (principle of explosion). In an axiomatic system, an axiom is called independent if it cannot be proven or disproven from other axioms in the system. A system is called independent if each of its underlying axioms is independent. Unlike consistency, independence is not a necessary requirement for a functioning axiomatic system — though it is usually sought after to minimize the number of axioms in the system. An axiomatic system is called complete if for every statement, either itself or its negation is derivable from the system's axioms (equivalently, every statement is capable of being proven true or false). Relative consistency Beyond consistency, relative consistency is also the mark of a worthwhile axiom system. This describes the scenario where the undefined terms of a first axiom system are provided definitions from a second, such that the axioms of the first are theorems of the second. A good example is the relative consistency of absolute geometry with respect to the theory of the real number system. Lines and points are undefined terms (also called primitive notions) in absolute geometry, but assigned meanings in the theory of real numbers in a way that is consistent with both axiom systems. Models A model for an axiomatic system is a well-defined set, which assigns meaning for the undefined terms presented in the system, in a manner that is correct with the relations defined in the system. The existence of a proves the consistency of a system. A model is called concrete if the meanings assigned are objects and relations from the real world, as opposed to an which is based on other axiomatic systems. Models can also be used to show the independence of an axiom in the system. By constructing a valid model for a subsystem without a specific axiom, we show that the omitted axiom is independent if its correctness does not necessarily follow from the subsystem. Two models are said to be isomorphic if a one-to-one correspondence can be found between their elements, in a manner that preserves their relationship. An axiomatic system for which every model is isomorphic to another is called (sometimes ). The property of categoriality (categoricity) ensures the completeness of a system, however the converse is not true: Completeness does not ensure the categoriality (categoricity) of a system, since two models can differ in properties that cannot be expressed by the semantics of the system. Example As an example, observe the following axiomatic system, based on first-order logic with additional semantics of the following countably infinitely many axioms added (these can be easily formalized as an axiom schema): (informally, there exist two different items). (informally, there exist three different items). Informally, this infinite set of axioms states that there are infinitely many different items. However, the concept of an infinite set cannot be defined within the system — let alone the cardinality of such a set. The system has at least two different models – one is the natural numbers (isomorphic to any other countably infinite set), and another is the real numbers (isomorphic to any other set with the cardinality of the continuum). In fact, it has an infinite number of models, one for each cardinality of an infinite set. However, the property distinguishing these models is their cardinality — a property which cannot be defined within the system. Thus the system is not categorial. However it can be shown to be complete, for example by using the Łoś–Vaught test. Axiomatic method Stating definitions and propositions in a way such that each new term can be formally eliminated by the priorly introduced terms requires primitive notions (axioms) to avoid infinite regress. This way of doing mathematics is called the axiomatic method. A common attitude towards the axiomatic method is logicism. In their book Principia Mathematica, Alfred North Whitehead and Bertrand Russell attempted to show that all mathematical theory could be reduced to some collection of axioms. More generally, the reduction of a body of propositions to a particular collection of axioms underlies the mathematician's research program. This was very prominent in the mathematics of the twentieth century, in particular in subjects based around homological algebra. The explication of the particular axioms used in a theory can help to clarify a suitable level of abstraction that the mathematician would like to work with. For example, mathematicians opted that rings need not be commutative, which differed from Emmy Noether's original formulation. Mathematicians decided to consider topological spaces more generally without the separation axiom which Felix Hausdorff originally formulated. The Zermelo–Fraenkel set theory, a result of the axiomatic method applied to set theory, allowed the "proper" formulation of set-theory problems and helped avoid the paradoxes of naïve set theory. One such problem was the continuum hypothesis. Zermelo–Fraenkel set theory, with the historically controversial axiom of choice included, is commonly abbreviated ZFC, where "C" stands for "choice". Many authors use ZF to refer to the axioms of Zermelo–Fraenkel set theory with the axiom of choice excluded. Today ZFC is the standard form of axiomatic set theory and as such is the most common foundation of mathematics. History Mathematical methods developed to some degree of sophistication in ancient Egypt, Babylon, India, and China, apparently without employing the axiomatic method. Euclid of Alexandria authored the earliest extant axiomatic presentation of Euclidean geometry and number theory. His idea begins with five undeniable geometric assumptions called axioms. Then, using these axioms, he established the truth of other propositions by proofs, hence the axiomatic method. Many axiomatic systems were developed in the nineteenth century, including non-Euclidean geometry, the foundations of real analysis, Cantor's set theory, Frege's work on foundations, and Hilbert's 'new' use of axiomatic method as a research tool. For example, group theory was first put on an axiomatic basis towards the end of that century. Once the axioms were clarified (that inverse elements should be required, for example), the subject could proceed autonomously, without reference to the transformation group origins of those studies. Issues Not every consistent body of propositions can be captured by a describable collection of axioms. In recursion theory, a collection of axioms is called recursive if a computer program can recognize whether a given proposition in the language is a theorem. Gödel's first incompleteness theorem then tells us that there are certain consistent bodies of propositions with no recursive axiomatization. Typically, the computer can recognize the axioms and logical rules for deriving theorems, and the computer can recognize whether a proof is valid, but to determine whether a proof exists for a statement is only soluble by "waiting" for the proof or disproof to be generated. The result is that one will not know which propositions are theorems and the axiomatic method breaks down. An example of such a body of propositions is the theory of the natural numbers, which is only partially axiomatized by the Peano axioms (described below). In practice, not every proof is traced back to the axioms. At times, it is not even clear which collection of axioms a proof appeals to. For example, a number-theoretic statement might be expressible in the language of arithmetic (i.e. the language of the Peano axioms) and a proof might be given that appeals to topology or complex analysis. It might not be immediately clear whether another proof can be found that derives itself solely from the Peano axioms. Any more-or-less arbitrarily chosen system of axioms is the basis of some mathematical theory, but such an arbitrary axiomatic system will not necessarily be free of contradictions, and even if it is, it is not likely to shed light on anything. Philosophers of mathematics sometimes assert that mathematicians choose axioms "arbitrarily", but it is possible that although they may appear arbitrary when viewed only from the point of view of the canons of deductive logic, that appearance is due to a limitation on the purposes that deductive logic serves. Example: The Peano axiomatization of natural numbers The mathematical system of natural numbers 0, 1, 2, 3, 4, ... is based on an axiomatic system first devised by the mathematician Giuseppe Peano in 1889. He chose the axioms, in the language of a single unary function symbol S (short for "successor"), for the set of natural numbers to be: There is a natural number 0. Every natural number a has a successor, denoted by Sa. There is no natural number whose successor is 0. Distinct natural numbers have distinct successors: if a ≠ b, then Sa ≠ Sb. If a property is possessed by 0 and also by the successor of every natural number it is possessed by, then it is possessed by all natural numbers ("Induction axiom"). Axiomatization In mathematics, axiomatization is the process of taking a body of knowledge and working backwards towards its axioms. It is the formulation of a system of statements (i.e. axioms) that relate a number of primitive terms — in order that a consistent body of propositions may be derived deductively from these statements. Thereafter, the proof of any proposition should be, in principle, traceable back to these axioms.	Mathematics	Axiomatic systems	null
188488	https://en.wikipedia.org/wiki/ZIP%20%28file%20format%29	ZIP (file format)	ZIP is an archive file format that supports lossless data compression. A ZIP file may contain one or more files or directories that may have been compressed. The ZIP file format permits a number of compression algorithms, though DEFLATE is the most common. This format was originally created in 1989 and was first implemented in PKWARE, Inc.'s PKZIP utility, as a replacement for the previous ARC compression format by Thom Henderson. The ZIP format was then quickly supported by many software utilities other than PKZIP. Microsoft has included built-in ZIP support (under the name "compressed folders") in versions of Microsoft Windows since 1998 via the "Plus! 98" addon for Windows 98. Native support was added as of the year 2000 in Windows ME. Apple has included built-in ZIP support in Mac OS X 10.3 (via BOMArchiveHelper, now Archive Utility) and later. Most free operating systems have built in support for ZIP in similar manners to Windows and macOS. ZIP files generally use the file extensions or and the MIME media type . ZIP is used as a base file format by many programs, usually under a different name. When navigating a file system via a user interface, graphical icons representing ZIP files often appear as a document or other object prominently featuring a zipper. History The file format was designed by Phil Katz of PKWARE and Gary Conway of Infinity Design Concepts. The format was created after Systems Enhancement Associates (SEA) filed a lawsuit against PKWARE claiming that the latter's archiving products, named PKARC, were derivatives of SEA's ARC archiving system. The name "zip" (meaning "move at high speed") was suggested by Katz's friend, Robert Mahoney. They wanted to imply that their product would be faster than ARC and other compression formats of the time. By distributing the zip file format within APPNOTE.TXT, compatibility with the zip file format proliferated widely on the public Internet during the 1990s. PKWARE and Infinity Design Concepts made a joint press release on February 14, 1989, releasing the file format into the public domain. Version history The .ZIP File Format Specification has its own version number, which does not necessarily correspond to the version numbers for the PKZIP tool, especially with PKZIP 6 or later. At various times, PKWARE has added preliminary features that allow PKZIP products to extract archives using advanced features, but PKZIP products that create such archives are not made available until the next major release. Other companies or organizations support the PKWARE specifications at their own pace. The .ZIP file format specification is formally named "APPNOTE - .ZIP File Format Specification" and it is published on the PKWARE.com website since the late 1990s. Several versions of the specification were not published. Specifications of some features such as BZIP2 compression, strong encryption specification and others were published by PKWARE a few years after their creation. The URL of the online specification was changed several times on the PKWARE website. A summary of key advances in various versions of the PKWARE specification: 2.0: (1993) File entries can be compressed with DEFLATE and use traditional PKWARE encryption (ZipCrypto). 2.1: (1996) Deflate64 compression 4.5: (2001) Documented 64-bit zip format. 4.6: (2001) BZIP2 compression (not published online until the publication of APPNOTE 5.2) 5.0: (2002) SES: DES, Triple DES, RC2, RC4 supported for encryption (not published online until the publication of APPNOTE 5.2) 5.2: (2003) AES encryption support for SES (defined in APPNOTE 5.1 that was not published online) and AES from WinZip ("AE-x"); corrected version of RC2-64 supported for SES encryption. 6.1: (2004) Documented certificate storage. 6.2.0: (2004) Documented Central Directory Encryption. 6.3.0: (2006) Documented Unicode (UTF-8) filename storage. Expanded list of supported compression algorithms (LZMA, PPMd+), encryption algorithms (Blowfish, Twofish), and hashes. 6.3.1: (2007) Corrected standard hash values for SHA-256/384/512. 6.3.2: (2007) Documented compression method 97 (WavPack). 6.3.3: (2012) Document formatting changes to facilitate referencing the PKWARE Application Note from other standards using methods such as the JTC 1 Referencing Explanatory Report (RER) as directed by JTC 1/SC 34 N 1621. 6.3.4: (2014) Updates the PKWARE, Inc. office address. 6.3.5: (2018) Documented compression methods 16, 96 and 99, DOS timestamp epoch and precision, added extra fields for keys and decryption, as well as typos and clarifications. 6.3.6: (2019) Corrected typographical error. 6.3.7: (2020) Added Zstandard compression method ID 20. 6.3.8: (2020) Moved Zstandard compression method ID from 20 to 93, deprecating the former. Documented method IDs 94 and 95 (MP3 and XZ respectively). 6.3.9: (2020) Corrected a typo in Data Stream Alignment description. 6.3.10: (2022) Added several z/OS attribute values for APPENDIX B. Added several additional 3rd party Extra Field mappings. WinZip, starting with version 12.1, uses the extension for ZIP files that use compression methods newer than DEFLATE; specifically, methods BZip, LZMA, PPMd, Jpeg and Wavpack. The last 2 are applied to appropriate file types when "Best method" compression is selected. Standardization In April 2010, ISO/IEC JTC 1 initiated a ballot to determine whether a project should be initiated to create an ISO/IEC International Standard format compatible with ZIP. The proposed project, entitled Document Packaging, envisaged a ZIP-compatible 'minimal compressed archive format' suitable for use with a number of existing standards including OpenDocument, Office Open XML and EPUB. In 2015, ISO/IEC 21320-1 "Document Container File — Part 1: Core" was published which states that "Document container files are conforming Zip files". It requires the following main restrictions of the ZIP file format: Files in ZIP archives may only be stored uncompressed, or using the "deflate" compression (i.e. compression method may contain the value "0" - stored or "8" - deflated). The encryption features are prohibited. The digital signature features (from SES) are prohibited. The "patched data" features (from PKPatchMaker) are prohibited. Archives may not span multiple volumes or be segmented. Design files are archives that store multiple files. ZIP allows contained files to be compressed using many different methods, as well as simply storing a file without compressing it. Each file is stored separately, allowing different files in the same archive to be compressed using different methods. Because the files in a ZIP archive are compressed individually, it is possible to extract them, or add new ones, without applying compression or decompression to the entire archive. This contrasts with the format of compressed tar files, for which such random-access processing is not easily possible. A directory is placed at the end of a ZIP file. This identifies what files are in the ZIP and identifies where in the ZIP that file is located. This allows ZIP readers to load the list of files without reading the entire ZIP archive. ZIP archives can also include extra data that is not related to the ZIP archive. This allows for a ZIP archive to be made into a self-extracting archive (application that decompresses its contained data), by prepending the program code to a ZIP archive and marking the file as executable. Storing the catalog at the end also makes possible to hide a zipped file by appending it to an innocuous file, such as a GIF image file. The format uses a 32-bit CRC algorithm and includes two copies of each entry metadata to provide greater protection against data loss. The CRC-32 algorithm was contributed by David Schwaderer and can be found in his book "C Programmers Guide to NetBIOS" published by Howard W. Sams & Co. Inc. Structure A ZIP file is correctly identified by the presence of an end of central directory record which is located at the end of the archive structure in order to allow the easy appending of new files. If the end of central directory record indicates a non-empty archive, the name of each file or directory within the archive should be specified in a central directory entry, along with other metadata about the entry, and an offset into the ZIP file, pointing to the actual entry data. This allows a file listing of the archive to be performed relatively quickly, as the entire archive does not have to be read to see the list of files. The entries within the ZIP file also include this information, for redundancy, in a local file header. Because ZIP files may be appended to, only files specified in the central directory at the end of the file are valid. Scanning a ZIP file for local file headers is invalid (except in the case of corrupted archives), as the central directory may declare that some files have been deleted and other files have been updated. For example, we may start with a ZIP file that contains files A, B and C. File B is then deleted and C updated. This may be achieved by just appending a new file C to the end of the original ZIP file and adding a new central directory that only lists file A and the new file C. When ZIP was first designed, transferring files by floppy disk was common, yet writing to disks was very time-consuming. If you had a large zip file, possibly spanning multiple disks, and only needed to update a few files, rather than reading and re-writing all the files, it would be substantially faster to just read the old central directory, append the new files then append an updated central directory. The order of the file entries in the central directory need not coincide with the order of file entries in the archive. Each entry stored in a ZIP archive is introduced by a local file header with information about the file such as the comment, file size and file name, followed by optional "extra" data fields, and then the possibly compressed, possibly encrypted file data. The "Extra" data fields are the key to the extensibility of the ZIP format. "Extra" fields are exploited to support the ZIP64 format, WinZip-compatible AES encryption, file attributes, and higher-resolution NTFS or Unix file timestamps. Other extensions are possible via the "Extra" field. ZIP tools are required by the specification to ignore Extra fields they do not recognize. The ZIP format uses specific 4-byte "signatures" to denote the various structures in the file. Each file entry is marked by a specific signature. The end of central directory record is indicated with its specific signature, and each entry in the central directory starts with the 4-byte central file header signature. There is no BOF or EOF marker in the ZIP specification. Conventionally the first thing in a ZIP file is a ZIP entry, which can be identified easily by its local file header signature. However, this is not necessarily the case, as this is not required by the ZIP specification - most notably, a self-extracting archive will begin with an executable file header. Tools that correctly read ZIP archives must scan for the end of central directory record signature, and then, as appropriate, the other, indicated, central directory records. They must not scan for entries from the top of the ZIP file, because (as previously mentioned in this section) only the central directory specifies where a file chunk starts and that it has not been deleted. Scanning could lead to false positives, as the format does not forbid other data to be between chunks, nor file data streams from containing such signatures. However, tools that attempt to recover data from damaged ZIP archives will most likely scan the archive for local file header signatures; this is made more difficult by the fact that the compressed size of a file chunk may be stored after the file chunk, making sequential processing difficult. Most of the signatures end with the short integer 0x4b50, which is stored in little-endian ordering. Viewed as an ASCII string this reads "PK", the initials of the inventor Phil Katz. Thus, when a ZIP file is viewed in a text editor the first two bytes of the file are usually "PK". (DOS, OS/2 and Windows self-extracting ZIPs have an EXE before the ZIP so start with "MZ"; self-extracting ZIPs for other operating systems may similarly be preceded by executable code for extracting the archive's content on that platform.) The specification also supports spreading archives across multiple file-system files. Originally intended for storage of large ZIP files across multiple floppy disks, this feature is now used for sending ZIP archives in parts over email, or over other transports or removable media. The FAT filesystem of DOS has a timestamp resolution of only two seconds; ZIP file records mimic this. As a result, the built-in timestamp resolution of files in a ZIP archive is only two seconds, though extra fields can be used to store more precise timestamps. The ZIP format has no notion of time zone, so timestamps are only meaningful if it is known what time zone they were created in. In September 2006, PKWARE released a revision of the ZIP specification providing for the storage of file names using UTF-8, finally adding Unicode compatibility to ZIP. File headers All multi-byte values in the header are stored in little-endian byte order. All length fields count the length in bytes. Local file header The extra field contains a variety of optional data such as OS-specific attributes. It is divided into records, each with at minimum a 16-bit signature and a 16-bit length. A ZIP64 local file extra field record, for example, has the signature 0x0001 and a length of 16 bytes (or more) so that two 64-bit values (the uncompressed and compressed sizes) may follow. Another common local file extension is 0x5455 (or "UT") which contains 32-bit UTC UNIX timestamps. This is immediately followed by the compressed data. Data descriptor If the bit at offset 3 (0x08) of the general-purpose flags field is set, then the CRC-32 and file sizes are not known when the header is written. If the archive is in Zip64 format, the compressed and uncompressed size fields are 8 bytes long instead of 4 bytes long (see section 4.3.9.2). The equivalent fields in the local header (or in the Zip64 extended information extra field in the case of archives in Zip64 format) are filled with zero, and the CRC-32 and size are appended in a 12-byte structure (optionally preceded by a 4-byte signature) immediately after the compressed data: Central directory file header (CDFH) The central directory file header entry is an expanded form of the local header: End of central directory record (EOCD) After all the central directory entries comes the end of central directory (EOCD) record, which marks the end of the ZIP file: This ordering allows a ZIP file to be created in one pass, but the central directory is also placed at the end of the file in order to facilitate easy removal of files from multiple-part (e.g. "multiple floppy-disk") archives, as previously discussed. Compression methods The .ZIP File Format Specification documents the following compression methods: Store (no compression), Shrink (LZW), Reduce (levels 1–4; LZ77 + probabilistic), Implode, Deflate, Deflate64, bzip2, LZMA, Zstandard, WavPack, PPMd, and a LZ77 variant provided by IBM z/OS CMPSC instruction. The most commonly used compression method is DEFLATE, which is described in IETF . Other methods mentioned, but not documented in detail in the specification include: PKWARE DCL Implode (old IBM TERSE), new IBM TERSE, IBM LZ77 z Architecture (PFS), and a JPEG variant. A "Tokenize" method was reserved for a third party, but support was never added. The word Implode is overused by PKWARE: the DCL/TERSE Implode is distinct from the old PKZIP Implode, a predecessor to Deflate. The DCL Implode is undocumented partially due to its proprietary nature held by IBM, but Mark Adler has nevertheless provided a decompressor called "blast" alongside zlib. Encryption ZIP supports a simple password-based symmetric encryption system generally known as ZipCrypto. It is documented in the ZIP specification, and known to be seriously flawed. In particular, it is vulnerable to known-plaintext attacks, which are in some cases made worse by poor implementations of random-number generators. Computers running under native Microsoft Windows without third-party archivers can open, but not create, ZIP files encrypted with ZipCrypto, but cannot extract the contents of files using different encryption. New features including new compression and encryption (e.g. AES) methods have been documented in the ZIP File Format Specification since version 5.2. A WinZip-developed AES-based open standard ("AE-x" in APPNOTE) is used also by 7-Zip and Xceed, but some vendors use other formats. PKWARE SecureZIP (SES, proprietary) also supports RC2, RC4, DES, Triple DES encryption methods, Digital Certificate-based encryption and authentication (X.509), and archive header encryption. It is, however, patented (see ). File name encryption is introduced in .ZIP File Format Specification 6.2, which encrypts metadata stored in Central Directory portion of an archive, but Local Header sections remain unencrypted. A compliant archiver can falsify the Local Header data when using Central Directory Encryption. As of version 6.2 of the specification, the Compression Method and Compressed Size fields within Local Header are not yet masked. ZIP64 The original format had a 4 GB (232 bytes) limit on various things (uncompressed size of a file, compressed size of a file, and total size of the archive), as well as a limit of 65,535 (216-1) entries in a ZIP archive. In version 4.5 of the specification (which is not the same as v4.5 of any particular tool), PKWARE introduced the "ZIP64" format extensions to get around these limitations, increasing the limits to 16 EB (264 bytes). In essence, it uses a "normal" central directory entry for a file, followed by an optional "zip64" directory entry, which has the larger fields. The format of the Local file header (LOC) and Central directory file header (CDFH) are the same in ZIP and ZIP64. However, ZIP64 specifies an extra field that may be added to those records at the discretion of the compressor, whose purpose is to store values that do not fit in the classic LOC or CDFH records. To signal that the actual values are stored in ZIP64 extra fields, they are set to 0xFFFF or 0xFFFFFFFF in the corresponding LOC or CDFH record. If one entry does not fit into the classic LOC or CDFH record, only that entry is required to be moved into a ZIP64 extra field. The other entries may stay in the classic record. Therefore, not all entries shown in the following table might be stored in a ZIP64 extra field. However, if they appear, their order must be as shown in the table. On the other hand, the format of EOCD for ZIP64 is slightly different from the normal ZIP version. It is also not necessarily the last record in the file. An End of Central Directory Locator follows (an additional 20 bytes at the end). The File Explorer in Windows XP does not support ZIP64, but the Explorer in Windows Vista and later do. Likewise, some extension libraries support ZIP64, such as DotNetZip, QuaZIP and IO::Compress::Zip in Perl. Python's built-in zipfile supports it since 2.5 and defaults to it since 3.4. OpenJDK's built-in java.util.zip supports ZIP64 from version Java 7. Android Java API support ZIP64 since Android 6.0. Mac OS Sierra's Archive Utility notably does not support ZIP64, and can create corrupt archives when ZIP64 would be required. However, the ditto command shipped with Mac OS will unzip ZIP64 files. More recent versions of Mac OS ship with info-zip's zip and unzip command line tools which do support Zip64: to verify run zip -v and look for "ZIP64_SUPPORT". Combination with other file formats The file format allows for a comment containing up to 65,535 (216−1) bytes of data to occur at the end of the file after the central directory. Also, because the central directory specifies the offset of each file in the archive with respect to the start, it is possible for the first file entry to start at an offset other than zero, although some tools might not process archive files that do not start with a file entry at offset zero. The program gzip, for example, happens to be able to extract an entry from a .ZIP file if it is at offset zero. This allows arbitrary data to occur in the file both before and after the ZIP archive data, and for the archive to still be read by a ZIP application. A side-effect of this is that it is possible to author a file that is both a working ZIP archive and another format, provided that the other format tolerates arbitrary data at its end, beginning, or middle. Self-extracting archives (SFX), of the form supported by WinZip, take advantage of this, in that they are executable () files that conform to the PKZIP AppNote.txt specification, and can be read by compliant zip tools or libraries. This property of the format, and of the JAR format which is a variant of ZIP, can be exploited to hide rogue content (such as harmful Java classes) inside a seemingly harmless file, such as a GIF image uploaded to the web. This so-called GIFAR exploit has been demonstrated as an effective attack against web applications such as Facebook. Limits The minimum size of a file is 22 bytes. Such an empty zip file contains only an End of Central Directory Record (EOCD): The maximum size for both the archive file and the individual files inside it is 4,294,967,295 bytes (232−1 bytes, or 4 GB minus 1 byte) for standard ZIP. For ZIP64, the maximum size is 18,446,744,073,709,551,615 bytes (264−1 bytes, or 16 EB minus 1 byte). Open extensions Seek-optimized (SOZip) profile A Seek-Optimized ZIP file (SOZip) profile has been proposed for the ZIP format. Such file contains one or several Deflate-compressed files that are organized and annotated such that a SOZip-aware reader can perform very fast random access (seek) within a compressed file. SOZip makes it possible to access large compressed files directly from a .zip file without prior decompression. It combines the use of ZLib block flushes issued at regular interval with a hidden index file mapping offsets of the uncompressed file to offsets in the compressed stream. ZIP readers that are not aware of that extension can read a SOZip-enabled file normally and ignore the extended features that support efficient seek capability. Proprietary extensions Extra field file format includes an extra field facility within file headers, which can be used to store extra data not defined by existing ZIP specifications, and which allow compliant archivers that do not recognize the fields to safely skip them. Header IDs 0–31 are reserved for use by PKWARE. The remaining IDs can be used by third-party vendors for proprietary usage. Strong encryption controversy When WinZip 9.0 public beta was released in 2003, WinZip introduced its own AES-256 encryption, using a different file format, along with the documentation for the new specification. The encryption standards themselves were not proprietary, but PKWARE had not updated APPNOTE.TXT to include Strong Encryption Specification (SES) since 2001, which had been used by PKZIP versions 5.0 and 6.0. WinZip technical consultant Kevin Kearney and StuffIt product manager Mathew Covington accused PKWARE of withholding SES, but PKZIP chief technology officer Jim Peterson claimed that certificate-based encryption was still incomplete. In another controversial move, PKWare applied for a patent on 16 July 2003 describing a method for combining ZIP and strong encryption to create a secure file. In the end, PKWARE and WinZip agreed to support each other's products. On 21 January 2004, PKWARE announced the support of WinZip-based AES compression format. In a later version of WinZip beta, it was able to support SES-based ZIP files. PKWARE eventually released version 5.2 of the .ZIP File Format Specification to the public, which documented SES. The Free Software project 7-Zip also supports AES, but not SES in ZIP files (as does its POSIX port p7zip). When using AES encryption under WinZip, the compression method is always set to 99, with the actual compression method stored in an AES extra data field. In contrast, Strong Encryption Specification stores the compression method in the basic file header segment of Local Header and Central Directory, unless Central Directory Encryption is used to mask/encrypt metadata. Implementation There are numerous tools available, and numerous libraries for various programming environments; licenses used include proprietary and free software. WinZip, WinRAR, Info-ZIP, ZipGenius, 7-Zip, PeaZip and B1 Free Archiver are well-known tools, available on various platforms. Some of those tools have library or programmatic interfaces. Some development libraries licensed under open source agreement are libzip, libarchive, and Info-ZIP. For Java: Java Platform, Standard Edition contains the package "java.util.zip" to handle standard files; the Zip64File library specifically supports large files (larger than 4 GB) and treats files using random access; and the Apache Ant tool contains a more complete implementation released under the Apache Software License. The Info-ZIP implementations of the format adds support for Unix filesystem features, such as user and group IDs, file permissions, and support for symbolic links. The Apache Ant implementation is aware of these to the extent that it can create files with predefined Unix permissions. The Info-ZIP implementations also know how to use the error correction capabilities built into the compression format. Some programs do not, and will fail on a file that has errors. The Info-ZIP Windows tools also support NTFS filesystem permissions, and will make an attempt to translate from NTFS permissions to Unix permissions or vice versa when extracting files. This can result in potentially unintended combinations, e.g. .exe files being created on NTFS volumes with executable permission denied. Versions of Microsoft Windows have included support for compression in Explorer since the Microsoft Plus! pack was released for Windows 98. Microsoft calls this feature "Compressed Folders". Not all features are supported by the Windows Compressed Folders capability. For example, encryption is not supported in Windows 10 Home edition, although it can decrypt. Unicode entry encoding is not supported until Windows 7, while split and spanned archives are not readable or writable by the Compressed Folders feature, nor is AES Encryption supported. Windows .zip support stemmed from an acquisition of "VisualZip" written by Dave Plummer. OpenDocument Format (ODF) started using the zip archive format in 2005, ODF is an open format for office documents of all types, this is the default file format used in Collabora Online, LibreOffice and others. Microsoft Office started using the zip archive format in 2006 for their Office Open XML .docx, .xlsx, .pptx, etc. files, which became the default file format with Microsoft Office 2007. Internationalization issues Versions of the format prior to 6.3.0 did not support storing file names in Unicode. According to the standard, file names should be stored in the CP437 encoding, which is standard for the IBM PC, but in practice, DOS archivers used the system's installed character encoding. The built-in archiver of Windows up to 11 also used the DOS encoding corresponding to the selected system language for backward compatibility when creating archives. Subsequently, the standard was updated to include two options for storing file names in Unicode: 1) when the 11th bit in the General purpose bit flag field is set, the file name in the "File name" field of the header should be considered as UTF-8 rather than a single-byte encoding, and 2) the Unicode Path Extra Field was added to store the file name in UTF-8 encoding. Some versions of archivers on the Windows platform have also used ANSI encoding in the past. Thus, to correctly extract files with names containing non-English characters, it is necessary: Check for the presence of the Unicode Path Extra Field, and if it exists, use the filename from it, encoded in UTF-8. Check for the presence of flag 11 in the General purpose bit flag field, and if it is set, consider the filename encoding in the "File name" field to be UTF-8. If the "packing OS" field contains the value 11 (NTFS, Windows), and the "version of the packer" field value is greater than or equal to 20, consider the filename encoding in the "File name" field to be the ANSI (Windows) encoding corresponding to the system locale if one can be determined; otherwise, use CP437. If the "packing OS" field contains the value 0 (FAT, DOS), and the "version of the packer" field value is between 25 and 40 inclusive, consider the filename encoding in the local header's "File name" field to be ANSI (Windows) encoding, and in the central header's "File name" field to be OEM (DOS) encoding, corresponding to the system locale if one can be determined; otherwise, use CP437. In other cases, if the "OS packing" field contains the value 0 (FAT, DOS), 6 (HPFS, OS/2), or 11 (NTFS, Windows), consider the filename encoding in the "File name" field to be OEM (DOS) encoding, corresponding to the system locale if one can be determined; otherwise, use CP437. In all other cases, consider the filename encoding in the "File name" field to be the system encoding of operating system unpacker is running on. Some implementations of zip unpackers did not implement this algorithm or only partially implemented it, as a result, when viewing the contents of an archive or extracting it, users saw a chaotic set of characters, known as "mojibake", instead of letters of the national alphabet. In 2016, this problem was solved in the far2l file and archive manager for Linux, BSD and Mac. In 2024, similar solution was added to the version of 7zip used in the Debian distribution and its derivatives, and to the version of unzip used in the Ubuntu distribution and its dervicatives. Legacy There are numerous other standards and formats using "zip" as part of their name. For example, zip is distinct from gzip, and the latter is defined in IETF . Both zip and gzip primarily use the DEFLATE algorithm for compression. Likewise, the ZLIB format (IETF ) also uses the DEFLATE compression algorithm, but specifies different headers for error and consistency checking. Other common, similarly named formats and programs with different native formats include 7-Zip, bzip2, and rzip. Concerns The theoretical maximum compression factor for a raw DEFLATE stream is about 1032 to one, but by exploiting the ZIP format in unintended ways, ZIP archives with compression ratios of billions to one can be constructed. These zip bombs unzip to extremely large sizes, overwhelming the capacity of the computer they are decompressed on.	Technology	File formats	null
188494	https://en.wikipedia.org/wiki/Clavicle	Clavicle	The clavicle, collarbone, or keybone is a slender, S-shaped long bone approximately long that serves as a strut between the shoulder blade and the sternum (breastbone). There are two clavicles, one on each side of the body. The clavicle is the only long bone in the body that lies horizontally. Together with the shoulder blade, it makes up the shoulder girdle. It is a palpable bone and, in people who have less fat in this region, the location of the bone is clearly visible. It receives its name from Latin clavicula 'little key' because the bone rotates along its axis like a key when the shoulder is abducted. The clavicle is the most commonly fractured bone. It can easily be fractured by impacts to the shoulder from the force of falling on outstretched arms or by a direct hit. Structure The collarbone is a thin doubly curved long bone that connects the arm to the trunk of the body. Located directly above the first rib, it acts as a strut to keep the scapula in place so that the arm can hang freely. At its rounded medial end (sternal end), it articulates with the manubrium of the sternum (breastbone) at the sternoclavicular joint. At its flattened lateral end (acromial end), it articulates with the acromion, a process of the scapula (shoulder blade), at the acromioclavicular joint. The rounded medial region (sternal region) of the shaft has a long curve laterally and anteriorly along two-thirds of the entire shaft. The flattened lateral region (acromial region) of the shaft has an even larger posterior curve to articulate with the acromion of the scapula. The medial region is the longest clavicular region as it takes up two-thirds of the entire shaft. The lateral region is both the widest clavicular region and thinnest clavicular region. The lateral end has a rough inferior surface that bears a ridge, the trapezoid line, and a slight rounded projection, the conoid tubercle (above the coracoid process). These surface features are attachment sites for muscles and ligaments of the shoulder. It can be divided into three parts: medial end, lateral end, and shaft. Medial end The medial end is also known as the sternal end. It is quadrangular and articulates with the clavicular notch of the manubrium of the sternum to form the sternoclavicular joint. The articular surface extends to the inferior aspect for articulation with the first costal cartilage. Lateral end The lateral end is also known as the acromial end. It is flat from above downward. It bears a facet that articulates with the shoulder to form the acromioclavicular joint. The area surrounding the joint gives an attachment to the joint capsule. The anterior border is concave forward and the posterior border is convex backward. Shaft The shaft is divided into two main regions, the medial region, and the lateral region. The medial region is also known as the sternal region, it is the longest clavicular region as it takes up two-thirds of the entire shaft. The lateral region is also known as the acromial region, it is both the widest clavicular region and thinnest clavicular region. Lateral region of the shaft The lateral region of the shaft has two borders and two surfaces. the anterior border is concave forward and gives origin to the deltoid muscle. the posterior border is convex and gives attachment to the trapezius muscle. the inferior surface has a ridge called the trapezoid line and a tubercle; the conoid tubercle for attachment with the trapezoid and the conoid ligament, part of the coracoclavicular ligament that serves to connect the collarbone with the coracoid process of the scapula. Development The collarbone is the first bone to begin the process of ossification (laying down of minerals onto a preformed matrix) during development of the embryo, during the fifth and sixth weeks of gestation. However, it is one of the last bones to finish ossification at about 21–25 years of age. Its lateral end is formed by intramembranous ossification while medially it is formed by endochondral ossification. It consists of a mass of cancellous bone surrounded by a compact bone shell. The cancellous bone forms via two ossification centres, one medial and one lateral, which fuse later on. The compact forms as the layer of fascia covering the bone stimulate the ossification of adjacent tissue. The resulting compact bone is known as a periosteal collar. The collarbone has a medullary cavity (marrow cavity) in its medial two-thirds. It is made up of spongy cancellous bone with a shell of compact bone. It is a dermal bone derived from elements originally attached to the skull. Variation The shape of the clavicle varies more than most other long bones. It is occasionally pierced by a branch of the supraclavicular nerve. In males the clavicle is usually longer and larger than in females. A study measuring 748 males and 252 females saw a difference in collarbone length between age groups 18–20 and 21–25 of about for males and females respectively. The left clavicle is usually longer and weaker than the right clavicle. The collarbones are sometimes partly or completely absent in cleidocranial dysostosis. The levator claviculae muscle, present in 2–3% of people, originates on the transverse processes of the upper cervical vertebrae and is inserted in the lateral half of the clavicle. Functions The collarbone serves several functions: It serves as a rigid support from which the scapula and free limb are suspended; an arrangement that keeps the upper limb away from the thorax so that the arm has maximum range of movement. Acting as a flexible, crane-like strut, it allows the scapula to move freely on the thoracic wall. Covering the cervicoaxillary canal, it protects the neurovascular bundle that supplies the upper limb. Transmits physical impacts from the upper limb to the axial skeleton. Muscle Muscles and ligaments that attach to the collarbone include: Clinical significance Acromioclavicular dislocation ("AC Separation") Degeneration of the clavicle Osteolysis Sternoclavicular dislocations A vertical line drawn from the mid-clavicle called the mid-clavicular line is used as a reference in describing cardiac apex beat during medical examination. It is also useful for evaluating an enlarged liver, and for locating the gallbladder which is between the mid-clavicular line and the transpyloric plane. Collarbone fracture Clavicle fractures (colloquially, a broken collarbone) occur as a result of injury or trauma. The most common type of fractures occur when a person falls horizontally on the shoulder or with an outstretched hand. A direct hit to the collarbone will also cause a break. In most cases, the direct hit occurs from the lateral side towards the medial side of the bone. The most common site of fracture is the junction between the two curvatures of the bone, which is the weakest point. This results in the sternocleidomastoid muscle lifting the medial aspect superiorly, which can result in perforation of the overlying skin. Other animals The clavicle first appears as part of the skeleton in primitive bony fish, where it is associated with the pectoral fin; they also have a bone called the cleithrum. In such fish, the paired clavicles run behind and below the gills on each side, and are joined by a solid symphysis on the fish's underside. They are, however, absent in cartilaginous fish and in the vast majority of living bony fish, including all of the teleosts. The earliest tetrapods retained this arrangement, with the addition of a diamond-shaped interclavicle between the base of the clavicles, although this is not found in living amphibians. The cleithrum disappeared early in the evolution of reptiles, and is not found in any living amniotes, but the interclavicle is present in most modern reptiles, and also in monotremes. In modern forms, however, there are a number of variations from the primitive pattern. For example, crocodilians and salamanders lack clavicles altogether (although crocodilians do retain the interclavicle), while in turtles, they form part of the armoured plastron. The interclavicle is absent in marsupials and placental mammals. In many mammals, the clavicles are also reduced, or even absent, to allow the scapula greater freedom of motion, which may be useful in fast-running animals. Though a number of fossil hominin (humans and chimpanzees) clavicles have been found, most of these are mere segments offering limited information on the form and function of the pectoral girdle. One exception is the clavicle of AL 333x6/9 attributed to Australopithecus afarensis which has a well-preserved sternal end. One interpretation of this specimen, based on the orientation of its lateral end and the position of the deltoid attachment area, suggests that this clavicle is distinct from those found in extant apes (including humans), and thus that the shape of the human shoulder dates back to less than . However, analyses of the clavicle in extant primates suggest that the low position of the scapula in humans is reflected mostly in the curvature of the medial portion of the clavicle rather than the lateral portion. This part of the bone is similar in A. afarensis and it is thus possible that this species had a high shoulder position similar to that in modern humans. In dinosaurs In dinosaurs, the main bones of the pectoral girdle were the scapula (shoulder blade) and the coracoid, both of which directly articulated with the clavicle. The clavicle was present in saurischian dinosaurs but largely absent in ornithischian dinosaurs. The place on the scapula where it articulated with the humerus (upper bone of the forelimb) is the called the glenoid. The clavicles fused in some theropod dinosaurs to form a furcula, which is the equivalent to a wishbone. In birds, the clavicles and interclavicle have fused to form a single Y-shaped bone, the furcula or "wishbone" which evolved from the clavicles found in coelurosaurian theropods. Additional media	Biology and health sciences	Skeletal system	Biology

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