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44312 | https://en.wikipedia.org/wiki/Invention | Invention | An invention is a unique or novel device, method, composition, idea or process. An invention may be an improvement upon a machine, product, or process for increasing efficiency or lowering cost. It may also be an entirely new concept. If an idea is unique enough either as a stand-alone invention or as a significant improvement over the work of others, it can be patented. A patent, if granted, gives the inventor a proprietary interest in the patent over a specific period of time, which can be licensed for financial gain.
An inventor creates or discovers an invention. The word inventor comes from the Latin verb invenire, invent-, to find. Although inventing is closely associated with science and engineering, inventors are not necessarily engineers or scientists. Due to advances in artificial intelligence, the term "inventor" no longer exclusively applies to an occupation (see human computers).
Some inventions can be patented. The system of patents was established to encourage inventors by granting limited-term, limited monopoly on inventions determined to be sufficiently novel, non-obvious, and useful. A patent legally protects the intellectual property rights of the inventor and legally recognizes that a claimed invention is actually an invention. The rules and requirements for patenting an invention vary by country and the process of obtaining a patent is often expensive.
Another meaning of invention is cultural invention, which is an innovative set of useful social behaviours adopted by people and passed on to others. The Institute for Social Inventions collected many such ideas in magazines and books. Invention is also an important component of artistic and design creativity. Inventions often extend the boundaries of human knowledge, experience or capability.
Types
Inventions are of three kinds: scientific-technological (including medicine), sociopolitical (including economics and law), and humanistic, or cultural.
Scientific-technological inventions include railroads, aviation, vaccination, hybridization, antibiotics, astronautics, holography, the atomic bomb, computing, the Internet, and the smartphone.
Sociopolitical inventions comprise new laws, institutions, and procedures that change modes of social behavior and establish new forms of human interaction and organization. Examples include the British Parliament, the US Constitution, the Manchester (UK) General Union of Trades, the Boy Scouts, the Red Cross, the Olympic Games, the United Nations, the European Union, and the Universal Declaration of Human Rights, as well as movements such as socialism, Zionism, suffragism, feminism, and animal-rights veganism.
Humanistic inventions encompass culture in its entirety and are as transformative and important as any in the sciences, although people tend to take them for granted. In the domain of linguistics, for example, many alphabets have been inventions, as are all neologisms (Shakespeare invented about 1,700 words). Literary inventions include the epic, tragedy, comedy, the novel, the sonnet, the Renaissance, neoclassicism, Romanticism, Symbolism, Aestheticism, Socialist Realism, Surrealism, postmodernism, and (according to Freud) psychoanalysis. Among the inventions of artists and musicians are oil painting, printmaking, photography, cinema, musical tonality, atonality, jazz, rock, opera, and the symphony orchestra. Philosophers have invented logic (several times), dialectics, idealism, materialism, utopia, anarchism, semiotics, phenomenology, behaviorism, positivism, pragmatism, and deconstruction. Religious thinkers are responsible for such inventions as monotheism, pantheism, Methodism, Mormonism, iconoclasm, puritanism, deism, secularism, ecumenism, and the Baháʼí Faith. Some of these disciplines, genres, and trends may seem to have existed eternally or to have emerged spontaneously of their own accord, but most of them have had inventors.
Process
Practical means
Ideas for an invention may be developed on paper or on a computer, by writing or drawing, by trial and error, by making models, by experimenting, by testing and/or by making the invention in its whole form. Brainstorming also can spark new ideas for an invention. Collaborative creative processes are frequently used by engineers, designers, architects and scientists. Co-inventors are frequently named on patents.
In addition, many inventors keep records of their working process – notebooks, photos, etc., including Leonardo da Vinci, Galileo Galilei, Evangelista Torricelli, Thomas Jefferson and Albert Einstein.
In the process of developing an invention, the initial idea may change. The invention may become simpler, more practical, it may expand, or it may even morph into something totally different. Working on one invention can lead to others too.
History shows that turning the concept of an invention into a working device is not always swift or direct. Inventions may also become more useful after time passes and other changes occur. For example, the parachute became more useful once powered flight was a reality.
Conceptual means
Invention is often a creative process. An open and curious mind allows an inventor to see beyond what is known. Seeing a new possibility, connection or relationship can spark an invention. Inventive thinking frequently involves combining concepts or elements from different realms that would not normally be put together. Sometimes inventors disregard the boundaries between distinctly separate territories or fields. Several concepts may be considered when thinking about invention.
Play
Play may lead to invention. Childhood curiosity, experimentation, and imagination can develop one's play instinct. Inventors feel the need to play with things that interest them, and to explore, and this internal drive brings about novel creations.
Sometimes inventions and ideas may seem to arise spontaneously while daydreaming, especially when the mind is free from its usual concerns. For example, both J. K. Rowling (the creator of Harry Potter) and Frank Hornby (the inventor of Meccano) first had their ideas while on train journeys.
In contrast, the successful aerospace engineer Max Munk advocated "aimful thinking".
Re-envisioning
To invent is to see anew. Inventors often envision a new idea, seeing it in their mind's eye. New ideas can arise when the conscious mind turns away from the subject or problem when the inventor's focus is on something else, or while relaxing or sleeping. A novel idea may come in a flash—a Eureka! moment. For example, after years of working to figure out the general theory of relativity, the solution came to Einstein suddenly in a dream "like a giant die making an indelible impress, a huge map of the universe outlined itself in one clear vision". Inventions can also be accidental, such as in the case of polytetrafluoroethylene (Teflon).
Insight
Insight can also be a vital element of invention. Such inventive insights may begin with questions, doubt or a hunch. It may begin by recognizing that something unusual or accidental may be useful or that it could open a new avenue for exploration. For example, the odd metallic color of plastic made by accidentally adding a thousand times too much catalyst led scientists to explore its metal-like properties, inventing electrically conductive plastic and light emitting plastic—an invention that won the Nobel Prize in 2000 and has led to innovative lighting, display screens, wallpaper and much more (see conductive polymer, and organic light-emitting diode or OLED).
Exploration
Invention is often an exploratory process with an uncertain or unknown outcome. There are failures as well as successes. Inspiration can start the process, but no matter how complete the initial idea, inventions typically must be developed.
Improvement
Inventors may, for example, try to improve something by making it more effective, healthier, faster, more efficient, easier to use, serve more purposes, longer lasting, cheaper, more ecologically friendly, or aesthetically different, lighter weight, more ergonomic, structurally different, with new light or color properties, etc.
Implementation
In economic theory, inventions are one of the chief examples of "positive externalities", a beneficial side effect that falls on those outside a transaction or activity. One of the central concepts of economics is that externalities should be internalized—unless some of the benefits of this positive externality can be captured by the parties, the parties are under-rewarded for their inventions, and systematic under-rewarding leads to under-investment in activities that lead to inventions. The patent system captures those positive externalities for the inventor or other patent owner so that the economy as a whole invests an optimum amount of resources in the invention process.
Comparison with innovation
In contrast to invention, innovation is the implementation of a creative idea that specifically leads to greater value or usefulness. That is, while an invention may be useless or have no value yet still be an invention, an innovation must have some sort of value, typically economic.
As defined by patent law
The term invention is also an important legal concept and central to patent law systems worldwide. As is often the case for legal concepts, its legal meaning is slightly different from common usage of the word. Additionally, the legal concept of invention is quite different in American and European patent law.
In Europe, the first test a patent application must pass is, "Is this an invention?" If it is, subsequent questions are whether it is new and sufficiently inventive. The implication—counter-intuitively—is that a legal invention is not inherently novel. Whether a patent application relates to an invention is governed by Article 52 of the European Patent Convention, that excludes, e.g., discoveries as such and software as such. The EPO Boards of Appeal decided that the technical character of an application is decisive for it to represent an invention, following an age-old Italian and German tradition. British courts do not agree with this interpretation. Following a 1959 Australian decision ("NRDC"), they believe that it is not possible to grasp the invention concept in a single rule. A British court once stated that the technical character test implies a "restatement of the problem in more imprecise terminology."
In the United States, all patent applications are considered inventions. The statute explicitly says that the American invention concept includes discoveries (35 USC § 100(a)), contrary to the European invention concept. The European invention concept corresponds to the American "patentable subject matter" concept: the first test a patent application is submitted to. While the statute (35 USC § 101) virtually poses no limits to patenting whatsoever, courts have decided in binding precedents that abstract ideas, natural phenomena and laws of nature are not patentable. Various attempts have been made to substantiate the "abstract idea" test, which suffers from abstractness itself, but none have succeeded. The last attempt so far was the "machine or transformation" test, but the U.S. Supreme Court decided in 2010 that it is merely an indication at best.
In India, invention means a new product or process that involves an inventive step, and capable of being made or used in an industry. Whereas, "new invention" means any invention that has not been anticipated in any prior art or used in the country or anywhere in the world.
In the arts
Invention has a long and important history in the arts. Inventive thinking has always played a vital role in the creative process. While some inventions in the arts are patentable, others are not because they cannot fulfill the strict requirements governments have established for granting them. (see patent).
Some inventions in art include the:
Collage and construction invented by Picasso
Readymade art invented by Marcel Duchamp
mobile invented by Alexander Calder
Combine invented by Robert Rauschenberg
Shaped painting invented by Frank Stella
Motion picture, the invention of which is attributed to Eadweard Muybridge
Video art invented by Nam June Paik
Likewise, Jackson Pollock invented an entirely new form of painting and a new kind of abstraction by dripping, pouring, splashing and splattering paint onto un-stretched canvas lying on the floor.
Inventive tools of the artist's trade also produced advances in creativity. Impressionist painting became possible because of newly invented collapsible, resealable metal paint tubes that facilitated spontaneous painting outdoors. Inventions originally created in the form of artwork can also develop other uses, e.g. Alexander Calder's mobile, which is now commonly used over babies' cribs. Funds generated from patents on inventions in art, design and architecture can support the realization of the invention or other creative work. Frédéric Auguste Bartholdi's 1879 design patent on the Statue of Liberty helped fund the famous statue because it covered small replicas, including those sold as souvenirs.
The timeline for invention in the arts lists the most notable artistic inventors.
Gender gap in inventions
Historically, women in many regions have been unrecognised for their inventive contributions (except Russia and France), despite being the sole inventor or co-inventor in inventions, including highly notable inventions. Notable examples include Margaret Knight who faced significant challenges in receiving credit for her inventions; Elizabeth Magie who was not credited for her invention of the game of Monopoly; and among other such examples, Chien-Shiung Wu whose male colleagues alone were awarded the Nobel Prize for their joint contributions to physics. Societal prejudice, institutional, educational and often legal patent barriers have both played a role in the gender invention gap. For example, although there could be found female patenters in US patent Office who also are likely to be helpful in their experience, still a patent applications made to the US Patent Office for inventions are less likely to succeed where the applicant have a "feminine" name, and additionally women could lose their independent legal patent rights to their husbands once married. | Technology | Technology: General | null |
44344 | https://en.wikipedia.org/wiki/Clover | Clover | Clovers, also called trefoils, are plants of the genus Trifolium (), consisting of about 300 species of flowering plants in the legume family Fabaceae originating in Europe. The genus has a cosmopolitan distribution with the highest diversity in the temperate Northern Hemisphere, but many species also occur in South America and Africa, including at high altitudes on mountains in the tropics.
They are small annual, biennial, or short-lived perennial herbaceous plants, typically growing up to tall. The leaves are trifoliate (rarely, they have less or more than three leaflets; the more leaflets the leaf has, the rarer it is; see four-leaf clover), with stipules adnate to the leaf-stalk, and heads or dense spikes of small red, purple, white, or yellow flowers; the small, few-seeded pods are enclosed in the calyx. Other closely related genera often called clovers include Melilotus (sweet clover) and Medicago (alfalfa or Calvary clover).
As legumes, clovers fix nitrogen using symbiotic bacteria in their root nodules, and are used as an alternative or supplement to synthetic nitrogen fertilizers. They are also valuable food source for grazing livestock and bees. The domestication of clover caused substantial increases in agricultural productivity.
Cultivation history
Clover was first domesticated in Spain in around the year 1000. During European urbanization, crop rotations involving clover became essential for replacing the fixed nitrogen exported to cities as food. Increased soil nitrogen levels from the spreading use of clover were one of the main reasons why European agricultural production in 1880 was about 275% of the production in 1750. Fields of clover, used as forage and newly-invented silage, became an important part of the rural landscape; adding clover made livestock feed more nutritious. Honey production also rose drastically, and clover remained the main nectar source for bees until the mid-twentieth century. Clover was carried around the world as a crop by European colonists, and some clover species became invasive in some areas.
Imports of guano and the development of the Haber-Bosch process in the 20th century substantially displaced clover as a crop, with negative effects on pollinators, but in the 1990s and 2010s, the cost of industrially-fixed nitrogen rose substantially, approximately doubling, and reviving interest in forage mixes that include clover. As the fixation process is energy-intensive, prices are closely tied to energy prices. The 21st century also saw interest in clover as a countermeasure to fight the global pollinator decline.
Cultivation
Several species of clover are extensively cultivated as fodder plants. The most widely cultivated clovers are white clover, Trifolium repens, and red clover, Trifolium pratense. Clover, either sown alone or in mixture with ryegrass, has for a long time formed a staple crop for silaging, for several reasons: it grows freely, shooting up again after repeated mowings; it produces an abundant crop; it is palatable to and nutritious for livestock; it fixes nitrogen using symbiotic bacteria in its root nodules, reducing the need for synthetic fertilizers; it grows in a great range of soils and climates; and it is appropriate for either pasturage or green composting.
In many areas, particularly on acidic soil, clover is short-lived because of a combination of insect pests, diseases and nutrient balance; this is known as "clover sickness". When crop rotations are managed so that clover does not recur at intervals shorter than eight years, it grows with much of its pristine vigor.
Clovers are most efficiently pollinated by bumblebees, which have declined as a result of agricultural intensification. Honeybees can also pollinate clover, and beekeepers are often in heavy demand from farmers with clover pastures. Farmers reap the benefits of increased reseeding that occurs with increased bee activity, which means that future clover yields remain abundant. Beekeepers benefit from the clover bloom, as clover is one of the main nectar sources for honeybees.
Trifolium repens, white or Dutch clover, is a perennial abundant in meadows and good pastures. The flowers are white or pinkish, becoming brown and deflexed as the corolla fades. Trifolium hybridum, alsike or Swedish clover, is a perennial which was introduced early in the 19th century and has now become naturalized in Britain. The flowers are white or rosy, and resemble those of Trifolium repens. Trifolium medium, meadow or zigzag clover, a perennial with straggling flexuous stems and rose-purple flowers, has potential for interbreeding with T. pratense to produce perennial crop plants.
Other species are: Trifolium arvense, hare's-foot trefoil; found in fields and dry pastures, a soft hairy plant with minute white or pale pink flowers and feathery sepals; Trifolium fragiferum, strawberry clover, with globose, rose-purple heads and swollen calyxes; Trifolium campestre, hop trefoil, on dry pastures and roadsides, the heads of pale yellow flowers suggesting miniature hops; and the somewhat similar Trifolium dubium, common in pastures and roadsides, with smaller heads and small yellow flowers turning dark brown.
Uses
Clover is foraged for by wildlife such as bears, game animals, and birds. Clover is edible by humans, although red clover should be avoided by pregnant women. The plant is a traditional Native American food, which is eaten both raw and after drying and smoking the roots. The seeds from the blossoms are used to make bread. It is also possible to make tea from the blossoms.
Symbolism
Shamrock, the traditional Irish symbol, which according to legend was coined by Saint Patrick for the Holy Trinity, is commonly associated with clover, although alternatively sometimes with the various species within the genus Oxalis, which are also trifoliate.
Clovers occasionally have four leaflets, instead of the usual three. These four-leaf clovers, like other rarities, are considered lucky. Clovers can also have five, six, or more leaflets, but these are rarer still. The clover's outer leaf structure varies in physical orientation.
The record for most leaflets is 63, set on August 2, 2023, by Yoshiharu Watanabe in Japan. The previous record holder, Shigeo Obara, had discovered an 18-leaf clover in 2002, a 21-leaf clover in 2008 and a 56-leaf clover in 2009, also in Japan.
A common idiom is "to be (or to live) in clover", meaning to live a carefree life of ease, comfort, or prosperity.
A common saying in surgery [regarding the appearance of wound after hemorrhoidectomy] is "If it looks like clover, the trouble is over; if it looks like dahlia, it’s surely a failure."
A cloverleaf interchange is named for the resemblance to the leaflets of a (four-leaf) clover when viewed from the air.
Phylogeny
The first extensive classification of Trifolium had been done by Michael Zohary and David Heller, and it was subsequently released in 1984. They divided the genus into eight sections: Lotoidea, Paramesus, Mistyllus, Vesicamridula, Chronosemium, Trifolium, Trichoecephalum, and Involucrarium, with Lotoidea placed most basally. Within this classification system, Trifolium repens falls within section Lotoidea, the largest and least heterogeneous section. Lotoidea contains species from America, Africa, and Eurasia, considered a clade because of their inflorescence shape, floral structure, and legume that protrudes from the calyx. However, these traits are not unique to the section, and are shared with many other species in other sections. Zohary and Heller argued that the presence of these traits in other sections proved the basal position of Lotoidea, because they were ancestral. Aside from considering this section basal, they did not propose relationships between other sections.
Since then, molecular data has both questioned and confirmed the proposed phylogeny from Zohary and Heller. A genus-wide molecular study has since proposed a new classification system, made up of two subgenera, Chronosemium and Trifolium. This recent reclassification further divides subgenus Trifolium into eight sections. The molecular data supports the monophyletic nature of three sections proposed by Zohary and Heller (Tripholium, Paramesus, and Trichoecepalum), but not of Lotoidea (members of this section have since been reclassified into five other sections). Other molecular studies, although smaller, support the need to reorganize Lotoidea.
Species
291 species of Trifolium are accepted:
Trifolium absconditum
Trifolium acaule Steud. ex A.Rich.
Trifolium affine C.Presl
Trifolium acutiflorum
Trifolium × adulterinum
Trifolium affine
Trifolium africanum Ser.
Trifolium aintabense Boiss. & Hausskn.
Trifolium albopurpureum Torr. & A. Gray – rancheria clover
Trifolium alexandrinum L. – Egyptian clover, berseem clover
Trifolium alpestre L. – owl-head clover
Trifolium alpinum L. – alpine clover
Trifolium alsadami
Trifolium amabile Kunth
Trifolium ambiguum M. Bieb.
Trifolium amoenum Greene – showy Indian clover
Trifolium amphianthum
Trifolium andersonii A. Gray – Anderson's clover or fiveleaf clover
Trifolium andinum Nutt. – Intermountain clover
Trifolium andricum Lassen
Trifolium angulatum Waldst. & Kit.
Trifolium angustifolium L.
Trifolium ankaratrense
Trifolium apertum Bobrov
Trifolium appendiculatum
Trifolium argutum Banks & Sol.
Trifolium arvense L. – hare's-foot clover
Trifolium attenuatum Greene
Trifolium aureum Pollich – large hop trefoil
Trifolium baccarinii Chiov.
Trifolium badium Schreb.
Trifolium barbigerum Torr. – bearded clover
Trifolium barbulatum
Trifolium barnebyi (Isely) Dorn & Lichvar
Trifolium batmanicum Katzn.
Trifolium beckwithii W.H.Brewer ex S.Watson – Beckwith's clover
Trifolium bejariense Moric.
Trifolium × bertrandii
Trifolium berytheum Boiss. & C.I.Blanche
Trifolium biebersteinii
Trifolium bifidum A.Gray – notchleaf clover
Trifolium bilineatum Fresen.
Trifolium billardierei Spreng.
Trifolium bithynicum
Trifolium bivonae Guss.
Trifolium blancheanum Boiss.
Trifolium bobrovii
Trifolium bocconei Savi
Trifolium boissieri Guss.
Trifolium bolanderi A.Gray
Trifolium bordsilovskyi
Trifolium brandegeei S.Watson
Trifolium breweri S. Watson – forest clover
Trifolium brutium Ten.
Trifolium buckwestiorum Isely – Santa Cruz clover
Trifolium bullatum Boiss. & Hausskn.
Trifolium burchellianum Ser.
Trifolium calcaricum J.L.Collins & Wieboldt
Trifolium calocephalum Fresen.
Trifolium campestre Schreb. – hop trefoil
Trifolium canescens Willd.
Trifolium carolinianum Michx.
Trifolium caudatum Boiss.
Trifolium cernuum Brot.
Trifolium cheranganiense J.B.Gillett
Trifolium cherleri L.
Trifolium chilaloense Thulin
Trifolium chilense Hook. & Arn.
Trifolium chlorotrichum Boiss. & Balansa
Trifolium ciliolatum Benth. – foothill clover
Trifolium circumdatum Kunze
Trifolium clusii Godr.
Trifolium clypeatum L.
Trifolium congestum Guss.
Trifolium constantinopolitanum Ser.
Trifolium cryptopodium Steud. ex A. Rich.
Trifolium cyathiferum Lindl. – cup clover
Trifolium dalmaticum Vis.
Trifolium dasyphyllum Torr. & A.Gray
Trifolium dasyurum C.Presl
Trifolium davisii E.Hossain
Trifolium decorum Chiov.
Trifolium dedeckerae
Trifolium depauperatum Desv. – cowbag clover, balloon sack clover, or poverty clover
Trifolium dichotomum Hook. & Arn.
Trifolium dichroanthoides Rech.f.
Trifolium dichroanthum Boiss.
Trifolium diffusum Ehrh.
Trifolium dolopium Heldr. & Hochst. ex Gibelli & Belli
Trifolium douglasii House
Trifolium dubium Sibth. – lesser hop trefoil
Trifolium echinatum M.Bieb.
Trifolium egrissicum
Trifolium elgonense J.B.Gillett
Trifolium elizabethiae
Trifolium eriocephalum Nutt. – woollyhead clover
Trifolium eriosphaerum Boiss.
Trifolium erubescens Fenzl
Trifolium euxinum Zohary
Trifolium eximium Stephan ex Ser.
Trifolium farayense
Trifolium fergan-karaeri
Trifolium fontanum
Trifolium fragiferum L. – strawberry clover
Trifolium friscanum (S.L.Welsh) S.L.Welsh
Trifolium fucatum Lindl. – bull clover or sour clover
Trifolium gemellum Pourr. ex Willd.
Trifolium gillettianum Jacq.-Fél.
Trifolium glanduliferum Boiss.
Trifolium globosum L.
Trifolium glomeratum L. – clustered clover or bush clover
Trifolium gordeievii (Kom.) Z.Wei
Trifolium gracilentum Torr. & A.Gray – pinpoint clover
Trifolium grandiflorum Schreb.
Trifolium gymnocarpon Nutt. – hollyleaf clover
Trifolium hatschbachii
Trifolium haussknechtii Boiss.
Trifolium haydenii Porter
Trifolium heldreichianum (Gibelli & Belli) Hausskn.
Trifolium hickeyi
Trifolium hirtum All. – rose clover
Trifolium howellii S.Watson – canyon clover or Howell's clover
Trifolium humile
Trifolium hybridum L. – Alsike clover
Trifolium hydrophilum
Trifolium incarnatum L. – crimson clover
Trifolium infamia-ponertii
Trifolium israeliticum Zohary & Katzn.
Trifolium isthmocarpum Brot.
Trifolium jokerstii Vincent & Rand.Morgan
Trifolium juliani Batt.
Trifolium kentuckiense Chapel & Vincent
Trifolium kingii S.Watson
Trifolium lanceolatum (J.B.Gillett) J.B.Gillett
Trifolium lappaceum L.
Trifolium latifolium (Hook.) Greene
Trifolium latinum Sebast.
Trifolium leibergii A.Nelson & J.F.Macbr. – Leiberg's clover
Trifolium lemmonii S.Watson – Lemmon's clover
Trifolium leucanthum M.Bieb.
Trifolium ligusticum Balb. ex Loisel.
Trifolium longidentatum Nábelek
Trifolium longipes Nutt. – longstalk clover
Trifolium lucanicum Gasp.
Trifolium lugardii Bullock
Trifolium lupinaster L.
Trifolium macilentum Greene
Trifolium macraei Hook. & Arn. – Chilean clover, double-head clover, or MacRae's clover
Trifolium macrocephalum (Pursh) Poir. – largehead clover
Trifolium masaiense J.B.Gillett
Trifolium mattirolianum Chiov.
Trifolium mazanderanicum Rech.f.
Trifolium medium L. – zigzag clover
Trifolium meduseum C.I.Blanche ex Boiss.
Trifolium meironense Zohary & Lerner
Trifolium mesogitanum
Trifolium michaelis
Trifolium michelianum Savi
Trifolium micranthum Viv.
Trifolium microcephalum Pursh – smallhead clover
Trifolium microdon Hook. & Arn. – thimble clover
Trifolium miegeanum Maire
Trifolium minutissimum
Trifolium modestum
Trifolium monanthum A.Gray – mountain carpet clover
Trifolium montanum L.
Trifolium multinerve A. Rich.
Trifolium multistriatum W.D.J.Koch
Trifolium mutabile Port.
Trifolium nanum Torr.
Trifolium nerimaniae M.Keskin
Trifolium × neyrautii
Trifolium nigrescens Viv.
Trifolium noricum Wulfen
Trifolium obscurum Savi
Trifolium obtusiflorum Hook. – clammy clover
Trifolium occidentale Coombe
Trifolium ochroleucon Huds. - sulphur clover
Trifolium oliganthum Steud. – fewflower clover
Trifolium olivaceum
Trifolium orbelicum
Trifolium ornithopodioides L.
Trifolium owyheense Gilkey
Trifolium pachycalyx Zohary
Trifolium palaestinum Boiss.
Trifolium pallescens Schreb.
Trifolium pallidum Waldst. & Kit.
Trifolium palmeri
Trifolium pamphylicum
Trifolium pannonicum Jacq. – Hungarian clover
Trifolium parnassi Boiss. & Spruner
Trifolium parryi A.Gray
Trifolium patens Schreb.
Trifolium patulum Tausch
Trifolium pauciflorum d'Urv.
Trifolium × permixtum
Trifolium peruvianum Vogel
Trifolium philistaeum Zohary
Trifolium phitosianum N.Böhling, Greuter & Raus
Trifolium phleoides Pourr. ex Willd.
Trifolium physanthum Hook. & Arn.
Trifolium physodes Steven ex M. Bieb.
Trifolium pichisermollii J.B.Gillett
Trifolium pignantii Fauché. & Chaub.
Trifolium pilczii Adamović
Trifolium pilulare Boiss.
Trifolium piorkowskii Rand.Morgan & A.L.Barber
Trifolium plebeium Boiss.
Trifolium plumosum Douglas ex Hook.
Trifolium polymorphum Poir.
Trifolium polyodon Greene
Trifolium polyphyllum C.A.Mey.
Trifolium polystachyum Fresen.
Trifolium praetermissum Greuter, Pleger & Raus.
Trifolium pratense L. – red clover
Trifolium productum
Trifolium prophetarum M. Hossain
Trifolium pseudomedium
Trifolium pseudostriatum Baker f.
Trifolium pulchellum
Trifolium purpureum Loisel.
Trifolium purseglovei J. B. Gillett
Trifolium quartinianum A. Rich.
Trifolium radicosum Boiss. & Hohen.
Trifolium rechingeri
Trifolium reflexum L. – buffalo clover
Trifolium repens L. – shamrock (white clover)
Trifolium resupinatum L. – Persian clover, shaftal
Trifolium retusum L.
Trifolium × retyezaticum
Trifolium rhizomatosum
Trifolium rhombeum
Trifolium riograndense Burkart
Trifolium rollinsii
Trifolium roussaeanum Boiss.
Trifolium rubens L.
Trifolium rueppellianum Fresen.
Trifolium salmoneum Mouterde
Trifolium sannineum
Trifolium sarosiense
Trifolium saxatile All.
Trifolium scabrum L.
Trifolium schimperi (Hochst.) A.Rich.
Trifolium schneideri
Trifolium × schwarzii
Trifolium scutatum Boiss.
Trifolium sebastiani Savi
Trifolium semipilosum Fresen.
Trifolium setiferum Boiss.
Trifolium simense Fresen.
Trifolium sintenisii Freyn
Trifolium siskiyouense J.M.Gillett
Trifolium somalense Taub. ex Harms
Trifolium sonorense
Trifolium spadiceum L.
Trifolium spananthum Thulin
Trifolium spumosum L.
Trifolium squamosum (or maritimum) L. – sea clover
Trifolium squarrosum L.
Trifolium stellatum L.
Trifolium steudneri Schweinf.
Trifolium stipulaceum Thunb.
Trifolium stoloniferum Muhl. ex A. Eaton – running buffalo clover
Trifolium stolzii Harms
Trifolium striatum L. – knotted clover
Trifolium strictum L.
Trifolium subterraneum L. – subterranean clover
Trifolium suffocatum L.
Trifolium sylvaticum Gérard
Trifolium tembense Fresen.
Trifolium tenuifolium
Trifolium thalii Vill.
Trifolium thompsonii C.V.Morton – Thompson's clover
Trifolium tomentosum L.
Trifolium × traplii
Trifolium triaristatum Bertero ex Colla
Trifolium trichocalyx A.Heller – Monterey clover
Trifolium trichocephalum M. Bieb.
Trifolium trichopterum Pančić
Trifolium tumens Steven ex M.Bieb.
Trifolium ukingense Harms
Trifolium uniflorum L.
Trifolium usambarense Taub.
Trifolium variegatum Nutt. – whitetip clover
Trifolium vavilovii Eig
Trifolium velebiticum Degen
Trifolium velenovskyi Vandas
Trifolium vernum Phil.
Trifolium vesiculosum Savi
Trifolium vestitum D.Heller & Zohary
Trifolium virginicum Small
Trifolium wentzelianum Harms
Trifolium wettsteinii Dörfl. & Hayek
Trifolium wigginsii J. M. Gillett
Trifolium willdenovii Spreng. − tomcat clover
Trifolium wormskioldii Lehm. – cow clover
Trifolium xanthinum
| Biology and health sciences | Fabales | Plants |
44346 | https://en.wikipedia.org/wiki/Shallot | Shallot | The shallot is a cultivar group of the onion. Until 2010, the (French red) shallot was classified as a separate species, Allium ascalonicum. The taxon was synonymized with Allium cepa (the common onion) in 2010, as the difference was too small to justify a separate species.
As part of the onion genus Allium, its close relatives include garlic, scallions, leeks, chives, and the Chinese onion.
Etymology and names
The names scallion and shallot are derived from the Old French eschalotte, by way of eschaloigne, from the Latin Ascalōnia caepa or Ascalonian onion, a Ascalōnia caepa or Ascalonian onion, a namesake of the ancient city of Ascalon.
The term shallot is usually applied to the French red shallot (Allium cepa var. aggregatum, or the A. cepa Aggregatum Group). It is also used for the Persian shallot or musir (A. stipitatum) from the Zagros Mountains in Iran and Iraq, and the French gray shallot (Allium oschaninii) which is also known as griselle or "true shallot"; it grows wild from Central to Southwest Asia. The name shallot is also used for a scallion in New South Wales, Australia and among English-speaking people in Quebec while the term French shallot refers to the plant referred to on this page. In most English-speaking nations, the name is pronounced with the emphasis on the last syllable in common with the French pronunciation, sha-lot, while the emphasis is commonly made on the first syllable, shall-ət, in the United States.
The term eschalot, derived from the French word échalote, can also be used to refer to the shallot.
Description and cultivation
Like garlic, shallots are formed in clusters of offsets with a head composed of multiple cloves. The skin colour of shallots can vary from golden brown to gray to rose red, and their off-white flesh is usually tinged with green or magenta.
Shallots are extensively cultivated for culinary uses, propagated by offsets. In some regions ("long-season areas"), the offsets are usually planted in autumn (September or October in the Northern Hemisphere). In some other regions, the suggested planting time for the principal crop is early spring (typically in February or the beginning of March in the Northern Hemisphere).
In planting, the tops of the bulbs should be kept a little above ground, and the soil surrounding the bulbs is often drawn away when the roots have taken hold. They come to maturity in summer, although fresh shallots can now be found year-round in supermarkets. Shallots should not be planted on ground recently manured. Shallots suffer damage from leek moth larvae, which mine into the leaves or bulbs of the plant.
Nutrition
A raw shallot is 80% water, 17% carbohydrates, 2.5% protein and contains negligible fat (table). In a reference amount of , raw shallot supplies 72 calories and is a rich source of vitamin B6 (27% of the Daily Value, DV), while providing moderate amounts of manganese (14% DV) and vitamin C (10% DV) (table). No other micronutrients are in significant content.
Uses
Culinary
Shallots are used in cooking. They may be pickled. Finely-sliced deep-fried shallots are used as a condiment in Asian cuisine, often served with porridge. Shallots taste similar to other cultivars of the common onion, but have a milder flavor. Like onions, when sliced, raw shallots release substances that irritate the human eye, resulting in production of tears.
Fresh shallots can be stored in a cool, dry area (0 to 4 °C, 32 to 40 °F, 60 to 70% RH) for six months or longer. Chopped, dried shallots are available.
Europe
In Europe, the Pikant, Atlas, and Ed's Red types of shallots are the most common.
Asia
Shallots are the traditional choice for many dishes in Sri Lankan cuisine, including pol sambola, lunu miris and many meat, fish and vegetable dishes.
In most Indian cuisines, the distinction between onions and shallots is weak; larger varieties of shallot are sometimes confused with small red onions and used interchangeably. Indeed, most parts of India use the regional name for onion interchangeably with shallot (Maharashtra, for instance, where both are called kanda). The southern regions of India distinguish shallots from onions in recipes more often, especially the much loved tiny varieties (about the width of a finger); these are widely used in curries and different types of sambar, a lentil-based dish. Shallots pickled in red vinegar are common in many Indian restaurants, served along with sauces and papad on the condiments tray. They are also used as a home remedy for sore throats, mixed with jaggery or sugar. In Nepal, shallots are used as one of the ingredients for making momo.
In Kashmir shallots are widely used in preparation of Wazwan Kashmiri cuisine, as they add distinct flavor and prevent curry from becoming black, which is common with onions.
In Iran shallots are used in various ways, the most common being grated shallot mixed into dense yogurt, a combination served in almost every restaurant when one orders grills or kebabs. Shallots are also used to make different types of torshi (ترشی), a sour Iranian side dish consisting of a variety of vegetables under vinegar, eaten with main dishes in small quantities. Shallot is also pickled—called shour (شور) in Persian—along with other vegetables to be served as torshi.
In Southeast Asian cuisines, such as those of Indonesia, Vietnam, Thailand, Cambodia, Malaysia, Philippines, Singapore and Brunei, both shallots and garlic are often used as elementary spices. Raw shallots can also accompany cucumbers when pickled in mild vinegar solution. They are also often chopped finely, then fried until golden brown, resulting in tiny crispy shallot chips called bawang goreng (fried shallots) in Indonesian, which can be bought ready-made from groceries and supermarkets. Shallots enhance the flavor of many Southeast Asian dishes, such as fried rice variants. They are also often present in noodle and slaw dishes. Crispy shallot chips are also used in southern Chinese cuisine. In Indonesia, shallots are sometimes pickled and added to several traditional foods; the pickles' sourness is thought to increase the appetite. In the southern Philippines, shallot bulbs and leaves are used to make the popular spicy Maranao condiment called palapa, which is used in the dish Piaparan.
The tubular green leaves of the plant can also be eaten and are very similar to the leaves of spring onions and chives.
Gallery
| Biology and health sciences | Root vegetables | Plants |
44363 | https://en.wikipedia.org/wiki/Wien%27s%20displacement%20law | Wien's displacement law | In physics, Wien's displacement law states that the black-body radiation curve for different temperatures will peak at different wavelengths that are inversely proportional to the temperature. The shift of that peak is a direct consequence of the Planck radiation law, which describes the spectral brightness or intensity of black-body radiation as a function of wavelength at any given temperature. However, it had been discovered by German physicist Wilhelm Wien several years before Max Planck developed that more general equation, and describes the entire shift of the spectrum of black-body radiation toward shorter wavelengths as temperature increases.
Formally, the wavelength version of Wien's displacement law states that the spectral radiance of black-body radiation per unit wavelength, peaks at the wavelength given by:
where is the absolute temperature and is a constant of proportionality called Wien's displacement constant, equal to or .
This is an inverse relationship between wavelength and temperature. So the higher the temperature, the shorter or smaller the wavelength of the thermal radiation. The lower the temperature, the longer or larger the wavelength of the thermal radiation. For visible radiation, hot objects emit bluer light than cool objects. If one is considering the peak of black body emission per unit frequency or per proportional bandwidth, one must use a different proportionality constant. However, the form of the law remains the same: the peak wavelength is inversely proportional to temperature, and the peak frequency is directly proportional to temperature.
There are other formulations of Wien's displacement law, which are parameterized relative to other quantities. For these alternate formulations, the form of the relationship is similar, but the proportionality constant, , differs.
Wien's displacement law may be referred to as "Wien's law", a term which is also used for the Wien approximation.
In "Wien's displacement law", the word displacement refers to how the intensity-wavelength graphs appear shifted (displaced) for different temperatures.
Examples
Wien's displacement law is relevant to some everyday experiences:
A piece of metal heated by a blow torch first becomes "red hot" as the very longest visible wavelengths appear red, then becomes more orange-red as the temperature is increased, and at very high temperatures would be described as "white hot" as shorter and shorter wavelengths come to predominate the black body emission spectrum. Before it had even reached the red hot temperature, the thermal emission was mainly at longer infrared wavelengths, which are not visible; nevertheless, that radiation could be felt as it warms one's nearby skin.
One easily observes changes in the color of an incandescent light bulb (which produces light through thermal radiation) as the temperature of its filament is varied by a light dimmer. As the light is dimmed and the filament temperature decreases, the distribution of color shifts toward longer wavelengths and the light appears redder, as well as dimmer.
A wood fire at 1500 K puts out peak radiation at about 2000 nanometers. 98% of its radiation is at wavelengths longer than 1000 nm, and only a tiny proportion at visible wavelengths (390–700 nanometers). Consequently, a campfire can keep one warm but is a poor source of visible light.
The effective temperature of the Sun is 5778 Kelvin. Using Wien's law, one finds a peak emission per nanometer (of wavelength) at a wavelength of about 500 nm, in the green portion of the spectrum near the peak sensitivity of the human eye. On the other hand, in terms of power per unit optical frequency, the Sun's peak emission is at 343 THz or a wavelength of 883 nm in the near infrared. In terms of power per percentage bandwidth, the peak is at about 635 nm, a red wavelength. About half of the Sun's radiation is at wavelengths shorter than 710 nm, about the limit of the human vision. Of that, about 12% is at wavelengths shorter than 400 nm, ultraviolet wavelengths, which is invisible to an unaided human eye. A large amount of the Sun's radiation falls in the fairly small visible spectrum and passes through the atmosphere.
The preponderance of emission in the visible range, however, is not the case in most stars. The hot supergiant Rigel emits 60% of its light in the ultraviolet, while the cool supergiant Betelgeuse emits 85% of its light at infrared wavelengths. With both stars prominent in the constellation of Orion, one can easily appreciate the color difference between the blue-white Rigel (T = 12100 K) and the red Betelgeuse (T ≈ 3800 K). While few stars are as hot as Rigel, stars cooler than the Sun or even as cool as Betelgeuse are very commonplace.
Mammals with a skin temperature of about 300 K emit peak radiation at around 10 μm in the far infrared. This is therefore the range of infrared wavelengths that pit viper snakes and passive IR cameras must sense.
When comparing the apparent color of lighting sources (including fluorescent lights, LED lighting, computer monitors, and photoflash), it is customary to cite the color temperature. Although the spectra of such lights are not accurately described by the black-body radiation curve, a color temperature (the correlated color temperature) is quoted for which black-body radiation would most closely match the subjective color of that source. For instance, the blue-white fluorescent light sometimes used in an office may have a color temperature of 6500 K, whereas the reddish tint of a dimmed incandescent light may have a color temperature (and an actual filament temperature) of 2000 K. Note that the informal description of the former (bluish) color as "cool" and the latter (reddish) as "warm" is exactly opposite the actual temperature change involved in black-body radiation.
Discovery
The law is named for Wilhelm Wien, who derived it in 1893 based on a thermodynamic argument. Wien considered adiabatic expansion of a cavity containing waves of light in thermal equilibrium. Using Doppler's principle, he showed that, under slow expansion or contraction, the energy of light reflecting off the walls changes in exactly the same way as the frequency. A general principle of thermodynamics is that a thermal equilibrium state, when expanded very slowly, stays in thermal equilibrium.
Wien himself deduced this law theoretically in 1893, following Boltzmann's thermodynamic reasoning. It had previously been observed, at least semi-quantitatively, by an American astronomer, Langley. This upward shift in with is familiar to everyone—when an iron is heated in a fire, the first visible radiation (at around 900 K) is deep red, the lowest frequency visible light. Further increase in causes the color to change to orange then yellow, and finally blue at very high temperatures (10,000 K or more) for which the peak in radiation intensity has moved beyond the visible into the ultraviolet.
The adiabatic principle allowed Wien to conclude that for each mode, the adiabatic invariant energy/frequency is only a function of the other adiabatic invariant, the frequency/temperature. From this, he derived the "strong version" of Wien's displacement law: the statement that the blackbody spectral radiance is proportional to for some function of a single variable. A modern variant of Wien's derivation can be found in the textbook by Wannier and in a paper by E. Buckingham
The consequence is that the shape of the black-body radiation function (which was not yet understood) would shift proportionally in frequency (or inversely proportionally in wavelength) with temperature. When Max Planck later formulated the correct black-body radiation function it did not explicitly include Wien's constant . Rather, the Planck constant was created and introduced into his new formula. From the Planck constant and the Boltzmann constant , Wien's constant can be obtained.
Peak differs according to parameterization
The results in the tables above summarize results from other sections of this article. Percentiles are percentiles of the Planck blackbody spectrum. Only 25 percent of the energy in the black-body spectrum is associated with wavelengths shorter than the value given by the peak-wavelength version of Wien's law.
Notice that for a given temperature, different parameterizations imply different maximal wavelengths. In particular, the curve of intensity per unit frequency peaks at a different wavelength than the curve of intensity per unit wavelength.
For example, using and parameterization by wavelength, the wavelength for maximal spectral radiance is with corresponding frequency . For the same temperature, but parameterizing by frequency, the frequency for maximal spectral radiance is with corresponding wavelength .
These functions are radiance density functions, which are probability density functions scaled to give units of radiance. The density function has different shapes for different parameterizations, depending on relative stretching or compression of the abscissa, which measures the change in probability density relative to a linear change in a given parameter. Since wavelength and frequency have a reciprocal relation, they represent significantly non-linear shifts in probability density relative to one another.
The total radiance is the integral of the distribution over all positive values, and that is invariant for a given temperature under any parameterization. Additionally, for a given temperature the radiance consisting of all photons between two wavelengths must be the same regardless of which distribution you use. That is to say, integrating the wavelength distribution from to will result in the same value as integrating the frequency distribution between the two frequencies that correspond to and , namely from to . However, the distribution shape depends on the parameterization, and for a different parameterization the distribution will typically have a different peak density, as these calculations demonstrate.
The important point of Wien's law, however, is that any such wavelength marker, including the median wavelength (or, alternatively, the wavelength below which any specified percentage of the emission occurs) is proportional to the reciprocal of temperature. That is, the shape of the distribution for a given parameterization scales with and translates according to temperature, and can be calculated once for a canonical temperature, then appropriately shifted and scaled to obtain the distribution for another temperature. This is a consequence of the strong statement of Wien's law.
Frequency-dependent formulation
For spectral flux considered per unit frequency (in hertz), Wien's displacement law describes a peak emission at the optical frequency given by:
or equivalently
where is a constant resulting from the maximization equation, is the Boltzmann constant, is the Planck constant, and is the absolute temperature. With the emission now considered per unit frequency, this peak now corresponds to a wavelength about 76% longer than the peak considered per unit wavelength. The relevant math is detailed in the next section.
Derivation from Planck's law
Parameterization by wavelength
Planck's law for the spectrum of black-body radiation predicts the Wien displacement law and may be used to numerically evaluate the constant relating temperature and the peak parameter value for any particular parameterization. Commonly a wavelength parameterization is used and in that case the black body spectral radiance (power per emitting area per solid angle) is:
Differentiating with respect to and setting the derivative equal to zero gives:
which can be simplified to give:
By defining:
the equation becomes one in the single variable x:
which is equivalent to:
This equation is solved by
where is the principal branch of the Lambert W function, and gives . Solving for the wavelength in millimetres, and using kelvins for the temperature yields:
Parameterization by frequency
Another common parameterization is by frequency. The derivation yielding peak parameter value is similar, but starts with the form of Planck's law as a function of frequency :
The preceding process using this equation yields:
The net result is:
This is similarly solved with the Lambert W function:
giving .
Solving for produces:
Parameterization by the logarithm of wavelength or frequency
Using the implicit equation yields the peak in the spectral radiance density function expressed in the parameter radiance per proportional bandwidth. (That is, the density of irradiance per frequency bandwidth proportional to the frequency itself, which can be calculated by considering infinitesimal intervals of (or equivalently ) rather of frequency itself.) This is perhaps a more intuitive way of presenting "wavelength of peak emission". That yields .
Mean photon energy as an alternate characterization
Another way of characterizing the radiance distribution is via the mean photon energy
where is the Riemann zeta function.
The wavelength corresponding to the mean photon energy is given by
Criticism
Marr and Wilkin (2012) contend that the widespread teaching of Wien's displacement law in introductory courses is undesirable, and it would be better replaced by alternate material. They argue that teaching the law is problematic because:
the Planck curve is too broad for the peak to stand out or be regarded as significant;
the location of the peak depends on the parameterization, and they cite several sources as concurring that "that the designation of any peak of the function is not meaningful and should, therefore, be de-emphasized";
the law is not used for determining temperatures in actual practice, direct use of the Planck function being relied upon instead.
They suggest that the average photon energy be presented in place of Wien's displacement law, as being a more physically meaningful indicator of changes that occur with changing temperature. In connection with this, they recommend that the average number of photons per second be discussed in connection with the Stefan–Boltzmann law. They recommend that the Planck spectrum be plotted as a "spectral energy density per fractional bandwidth distribution," using a logarithmic scale for the wavelength or frequency.
| Physical sciences | Thermodynamics | Physics |
44364 | https://en.wikipedia.org/wiki/Black%20body | Black body | A black body or blackbody is an idealized physical body that absorbs all incident electromagnetic radiation, regardless of frequency or angle of incidence. The radiation emitted by a black body in thermal equilibrium with its environment is called black-body radiation. The name "black body" is given because it absorbs all colors of light. In contrast, a white body is one with a "rough surface that reflects all incident rays completely and uniformly in all directions."
A black body in thermal equilibrium (that is, at a constant temperature) emits electromagnetic black-body radiation. The radiation is emitted according to Planck's law, meaning that it has a spectrum that is determined by the temperature alone (see figure at right), not by the body's shape or composition.
An ideal black body in thermal equilibrium has two main properties:
It is an ideal emitter: at every frequency, it emits as much or more thermal radiative energy as any other body at the same temperature.
It is a diffuse emitter: measured per unit area perpendicular to the direction, the energy is radiated isotropically, independent of direction.
Real materials emit energy at a fraction—called the emissivity—of black-body energy levels. By definition, a black body in thermal equilibrium has an emissivity . A source with a lower emissivity, independent of frequency, is often referred to as a gray body.
Constructing black bodies with an emissivity as close to 1 as possible remains a topic of current interest.
In astronomy, the radiation from stars and planets is sometimes characterized in terms of an effective temperature, the temperature of a black body that would emit the same total flux of electromagnetic energy.
Definition
Isaac Newton introduced the notion of a black body in his 1704 book Opticks, with query 6 of the book stating:The idea of a black body originally was introduced by Gustav Kirchhoff in 1860 as follows:
A more modern definition drops the reference to "infinitely small thicknesses":
Idealizations
This section describes some concepts developed in connection with black bodies.
Cavity with a hole
A widely used model of a black surface is a small hole in a cavity with walls that are opaque to radiation. Radiation incident on the hole will pass into the cavity, and is very unlikely to be re-emitted if the cavity is large. Lack of any re-emission, means that the hole is behaving like a perfect black surface. The hole is not quite a perfect black surface—in particular, if the wavelength of the incident radiation is greater than the diameter of the hole, part will be reflected. Similarly, even in perfect thermal equilibrium, the radiation inside a finite-sized cavity will not have an ideal Planck spectrum for wavelengths comparable to or larger than the size of the cavity.
Suppose the cavity is held at a fixed temperature T and the radiation trapped inside the enclosure is at thermal equilibrium with the enclosure. The hole in the enclosure will allow some radiation to escape. If the hole is small, radiation passing in and out of the hole has negligible effect upon the equilibrium of the radiation inside the cavity. This escaping radiation will approximate black-body radiation that exhibits a distribution in energy characteristic of the temperature T and does not depend upon the properties of the cavity or the hole, at least for wavelengths smaller than the size of the hole. See the figure in the Introduction for the spectrum as a function of the frequency of the radiation, which is related to the energy of the radiation by the equation E = hf, with E = energy, h = Planck constant, f = frequency.
At any given time the radiation in the cavity may not be in thermal equilibrium, but the second law of thermodynamics states that if left undisturbed it will eventually reach equilibrium, although the time it takes to do so may be very long. Typically, equilibrium is reached by continual absorption and emission of radiation by material in the cavity or its walls. Radiation entering the cavity will be "thermalized" by this mechanism: the energy will be redistributed until the ensemble of photons achieves a Planck distribution. The time taken for thermalization is much faster with condensed matter present than with rarefied matter such as a dilute gas. At temperatures below billions of Kelvin, direct photon–photon interactions are usually negligible compared to interactions with matter. Photons are an example of an interacting boson gas, and as described by the H-theorem, under very general conditions any interacting boson gas will approach thermal equilibrium.
Transmission, absorption, and reflection
A body's behavior with regard to thermal radiation is characterized by its transmission τ, absorption α, and reflection ρ.
The boundary of a body forms an interface with its surroundings, and this interface may be rough or smooth. A nonreflecting interface separating regions with different refractive indices must be rough, because the laws of reflection and refraction governed by the Fresnel equations for a smooth interface require a reflected ray when the refractive indices of the material and its surroundings differ. A few idealized types of behavior are given particular names:
An opaque body is one that transmits none of the radiation that reaches it, although some may be reflected. That is, τ = 0 and α + ρ = 1.
A transparent body is one that transmits all the radiation that reaches it. That is, τ = 1 and α = ρ = 0.
A grey body is one where α, ρ and τ are constant for all wavelengths; this term also is used to mean a body for which α is temperature- and wavelength-independent.
A white body is one for which all incident radiation is reflected uniformly in all directions: τ = 0, α = 0, and ρ = 1.
For a black body, τ = 0, α = 1, and ρ = 0. Planck offers a theoretical model for perfectly black bodies, which he noted do not exist in nature: besides their opaque interior, they have interfaces that are perfectly transmitting and non-reflective.
Kirchhoff's perfect black bodies
Kirchhoff in 1860 introduced the theoretical concept of a perfect black body with a completely absorbing surface layer of infinitely small thickness, but Planck noted some severe restrictions upon this idea. Planck noted three requirements upon a black body: the body must (i) allow radiation to enter but not reflect; (ii) possess a minimum thickness adequate to absorb the incident radiation and prevent its re-emission; (iii) satisfy severe limitations upon scattering to prevent radiation from entering and bouncing back out. As a consequence, Kirchhoff's perfect black bodies that absorb all the radiation that falls on them cannot be realized in an infinitely thin surface layer, and impose conditions upon scattering of the light within the black body that are difficult to satisfy.
Realizations
A realization of a black body refers to a real world, physical embodiment. Here are a few.
Cavity with a hole
In 1898, Otto Lummer and Ferdinand Kurlbaum published an account of their cavity radiation source. Their design has been used largely unchanged for radiation measurements to the present day. It was a hole in the wall of a platinum box, divided by diaphragms, with its interior blackened with iron oxide. It was an important ingredient for the progressively improved measurements that led to the discovery of Planck's law. A version described in 1901 had its interior blackened with a mixture of chromium, nickel, and cobalt oxides. | Physical sciences | Thermodynamics | Physics |
44365 | https://en.wikipedia.org/wiki/Typewriter | Typewriter | A typewriter is a mechanical or electromechanical machine for typing characters. Typically, a typewriter has an array of keys, and each one causes a different single character to be produced on paper by striking an inked ribbon selectively against the paper with a type element. Thereby, the machine produces a legible written document composed of ink and paper. By the end of the 19th century, a person who used such a device was also referred to as a typewriter.
The first commercial typewriters were introduced in 1874, but did not become common in offices in the United States until after the mid-1880s. The typewriter quickly became an indispensable tool for practically all writing other than personal handwritten correspondence. It was widely used by professional writers, in offices, in business correspondence in private homes, and by students preparing written assignments.
Typewriters were a standard fixture in most offices up to the 1980s. After that, they began to be largely supplanted by personal computers running word processing software. Nevertheless, typewriters remain common in some parts of the world. For example, typewriters are still used in many Indian cities and towns, especially in roadside and legal offices, due to a lack of continuous, reliable electricity.
The QWERTY keyboard layout, developed for typewriters in the 1870s, remains the de facto standard for English-language computer keyboards. The origins of this layout still need to be clarified. Similar typewriter keyboards, with layouts optimised for other languages and orthographies, emerged soon afterward, and their layouts have also become standard for computer keyboards in their respective markets.
History
Although many modern typewriters have one of several similar designs, their invention was incremental, developed by numerous inventors working independently or in competition with each other over a series of decades. As with the automobile, the telephone, and telegraph, several people contributed insights and inventions that eventually resulted in ever more commercially successful instruments. Historians have estimated that some form of the typewriter was invented 52 times as thinkers and tinkerers tried to come up with a workable design.
Some early typing instruments include:
In 1575, an Italian printmaker, Francesco Rampazetto, invented the , a machine to impress letters in papers.
In 1714, Henry Mill obtained a patent in Britain for a machine that, from the patent, appears to have been similar to a typewriter. The patent shows that this machine was created: "[he] hath by his great study and paines & expence invented and brought to perfection an artificial machine or method for impressing or transcribing of letters, one after another, as in writing, whereby all writing whatsoever may be engrossed in paper or parchment so neat and exact as not to be distinguished from print; that the said machine or method may be of great use in settlements and public records, the impression being deeper and more lasting than any other writing, and not to be erased or counterfeited without manifest discovery."
In 1802, Italian Agostino Fantoni developed a particular typewriter to enable his blind sister to write.
Between 1801 and 1808, Italian Pellegrino Turri invented a typewriter for his blind friend Countess Carolina Fantoni da Fivizzano.
In 1823, Italian Pietro Conti da Cilavegna invented a new model of the typewriter, the , also known as .
In 1829, American William Austin Burt patented a machine called the "Typographer" which, in common with many other early machines, is listed as the "first typewriter". The London Science Museum describes it merely as "the first writing mechanism whose invention was documented", but even that claim may be excessive since Turri's invention pre-dates it.
By the mid-19th century, the increasing pace of business communication had created a need to mechanize the writing process. Stenographers and telegraphers could take down information at rates up to 130 words per minute, whereas a writer with a pen was limited to a maximum of 30 words per minute (the 1853 speed record).
From 1829 to 1870, many printing or typing machines were patented by inventors in Europe and America, but none went into commercial production.
American Charles Thurber developed multiple patents, of which his first in 1843 was created as an aid to blind people, such as the 1845 Chirographer.
In 1855, the Italian Giuseppe Ravizza created a prototype typewriter called Cembalo scrivano o macchina da scrivere a tasti ("Scribe harpsichord, or machine for writing with keys"). It was an advanced machine that let the user see the writing as it was typed.
In 1861, Father Francisco João de Azevedo, a Brazilian priest, made his typewriter with basic materials and tools, such as wood and knives. In that same year, the Brazilian emperor D. Pedro II, presented a gold medal to Father Azevedo for this invention. Many Brazilian people, as well as the Brazilian federal government recognize Fr. Azevedo as the inventor of the typewriter, a claim that has been the subject of some controversy.
In 1865, John Pratt, of Centre, Alabama (US), built a machine called the Pterotype which appeared in an 1867 Scientific American article and inspired other inventors.
Between 1864 and 1867, , a carpenter from South Tyrol (then part of Austria) developed several models and a fully functioning prototype typewriter in 1867.
Hansen Writing Ball
In 1865, Rev. Rasmus Malling-Hansen of Denmark invented the Hansen Writing Ball, which went into commercial production in 1870 and was the first commercially sold typewriter. It was a success in Europe and was reported as being used in offices on the European continent as late as 1909.
Malling-Hansen used a solenoid escapement to return the carriage on some of his models, which makes him a candidate for the title of inventor of the first "electric" typewriter.
The Hansen Writing Ball was produced with only upper-case characters. The Writing Ball was a template for inventor Frank Haven Hall to create a derivative that would produce letter prints cheaper and faster.
Malling-Hansen developed his typewriter further through the 1870s and 1880s and made many improvements, but the writing head remained the same. On the first model of the writing ball from 1870, the paper was attached to a cylinder inside a wooden box. In 1874, the cylinder was replaced by a carriage, moving beneath the writing head. Then, in 1875, the well-known "tall model" was patented, which was the first of the writing balls that worked without electricity. Malling-Hansen attended the world exhibitions in Vienna in 1873 and Paris in 1878 and he received the first-prize for his invention at both exhibitions.
Sholes and Glidden typewriter
The first typewriter to be commercially successful was patented in 1868 by Americans Christopher Latham Sholes, Frank Haven Hall, Carlos Glidden and Samuel W. Soule in Milwaukee, Wisconsin, although Sholes soon disowned the machine and refused to use or even recommend it. The working prototype was made by clock-maker and machinist Matthias Schwalbach. Hall, Glidden and Soule sold their shares in the patent (US 79,265) to Sholes and James Densmore, who made an agreement with E. Remington and Sons (then famous as a manufacturer of sewing machines) to commercialize the machine as the Sholes and Glidden Type-Writer. This was the origin of the term typewriter.
Remington began production of its first typewriter on March 1, 1873, in Ilion, New York. It had a QWERTY keyboard layout, which, because of the machine's success, was slowly adopted by other typewriter manufacturers. As with most other early typewriters, because the typebars strike upwards, the typist could not see the characters as they were typed.
Index typewriter
The index typewriter came into the market in the early 1880s. The index typewriter uses a pointer or stylus to choose a letter from an index. The pointer is mechanically linked so that the letter chosen could then be printed, most often by the activation of a lever.
The index typewriter was briefly popular in niche markets. Although they were slower than keyboard type machines, they were mechanically simpler and lighter. They were therefore marketed as being suitable for travellers and, because they could be produced more cheaply than keyboard machines, as budget machines for users who needed to produce small quantities of typed correspondence. For example, the Simplex Typewriter Company made index typewriters for 1/40 the price of a Remington typewriter.
The index typewriter's niche appeal however soon disappeared as, on the one hand new keyboard typewriters became lighter and more portable, and on the other refurbished second-hand machines began to become available. The last widely available western index machine was the Mignon typewriter produced by AEG which was produced until 1934. Considered one of the very best of the index typewriters, part of the Mignon's popularity was that it featured interchangeable indexes as well as type, fonts and character sets. This is something very few keyboard machines were capable of--and only at considerable added cost.
Although they were pushed out of the market in most of the world by keyboard machines, successful Japanese and Chinese typewriters typewriters are of the index type--albeit with a very much larger index and number of type elements.
Embossing tape label makers are the most common index typewriters today, and perhaps the most common typewriters of any kind still being manufactured.
The platen was mounted on a carriage that moved horizontally to the left, automatically advancing the typing position, after each character was typed. The carriage-return lever at the far left was then pressed to the right to return the carriage to its starting position and rotating the platen to advance the paper vertically. A small bell was struck a few characters before the right hand margin was reached to warn the operator to complete the word and then use the carriage-return lever.
Other typewriters
1884 – Hammond "Ideal" typewriter with case, by Hammond Typewriter Company Limited, United States. Despite an unusual, curved keyboard (see picture in citation), the Hammond became popular because of its superior print quality and changeable typeface. Invented by James Hammond of Boston, Massachusetts in 1880, and commercially released in 1884. The type is carried on a pair of interchangeable rotating sectors, one controlled by each half of the keyboard. A small hammer pushes the paper against the ribbon and type sector to print each character. The mechanism was later adapted to give a straight QWERTY keyboard and proportional spacing.
1888 – Fitch typewriter – Made by the Fitch Typewriter Company, Brooklyn, N.Y. and later in the UK with a slightly different look. Operators of the early typewriters had to work "blind": the typed text emerged only after several lines had been completed or the carriage was lifted to look underneath at the page. The Fitch was one of the first machines to allow prompt correction of mistakes with its visible writing; it was said to be the second machine operating on the visible writing system. The typebars were positioned behind the paper and the writing area faced upwards so that the result could be seen instantly. A curved frame kept the emerging paper from obscuring the keyboard, but the Fitch was soon eclipsed by machines in which the paper could be fed more conveniently at the rear.
1893 – Gardner typewriter. This typewriter, patented by Mr J Gardner in 1893, was an attempt to reduce the size and cost. Although it prints 84 symbols, it has only 14 keys and two change-case keys. Several characters are indicated on each key and the character printed is determined by the position of the case keys, which choose one of six cases.
1896 – The "Underwood 1 typewriter, 10" Pica, No. 990". This was the first typewriter with a typing area fully visible to the typist until a key is struck. These features, copied by all subsequent typewriters, allowed the typist to see and if necessary correct the typing as it proceeded. The mechanism was developed in the US by Franz X. Wagner from about 1892 and taken up, in 1895, by John T. Underwood (1857–1937), a producer of office supplies.
Standardization
By about 1910, the "manual" or "mechanical" typewriter had reached a somewhat standardized design. There were minor variations from one manufacturer to another, but most typewriters followed the concept that each key was attached to a typebar that had the corresponding letter molded, in reverse, into its striking head. When a key was struck briskly and firmly, the typebar hit a ribbon (usually made of inked fabric), making a printed mark on the paper wrapped around a cylindrical platen.
The platen was mounted on a carriage that moved horizontally to the left, automatically advancing the typing position, after each character was typed. The carriage-return lever at the far left was then pressed to the right to return the carriage to its starting position and rotating the platen to advance the paper vertically. A small bell was struck a few characters before the right hand margin was reached to warn the operator to complete the word and then use the carriage-return lever. Typewriters for languages written right-to-left operate in the opposite direction.
By 1900, notable typewriter manufacturers included E. Remington and Sons, IBM, Godrej, Imperial Typewriter Company, Oliver Typewriter Company, Olivetti, Royal Typewriter Company, Smith Corona, Underwood Typewriter Company, Facit, Adler, and Olympia-Werke.
After the market had matured under the market dominance of large companies from Britain, Europe and the United States—but before the advent of daisywheel and electronic machines—the typewriter market faced strong competition from less expensive typewriters from Asia, including Brother Industries and Silver Seiko Ltd. of Japan.
Frontstriking
In most of the early typewriters, the typebars struck upward against the paper, pressed against the bottom of the platen, so the typist could not see the text as it was typed. What was typed was not visible until a carriage return caused it to scroll into view.
The difficulty with any other arrangement was ensuring the typebars fell back into place reliably when the key was released. This was eventually achieved with various ingenious mechanical designs and so-called "visible typewriters" which used frontstriking, in which the typebars struck forward against the front side of the platen, became standard. One of the first was the Daugherty Visible, introduced in 1893, which also introduced the four-bank keyboard that became standard, although the Underwood which came out two years later was the first major typewriter with these features.
Shift key
A significant innovation was the shift key, introduced with the Remington No. 2 in 1878. This key physically "shifted" either the basket of typebars, in which case the typewriter is described as "basket shift", or the paper-holding carriage, in which case the typewriter is described as "carriage shift". Either mechanism caused a different portion of the typebar to come in contact with the ribbon/platen.
The result is that each typebar could type two different characters, cutting the number of keys and typebars in half (and simplifying the internal mechanisms considerably). The obvious use for this was to allow letter keys to type both upper and lower case, but normally the number keys were also duplexed, allowing access to special symbols such as percent, , and ampersand, .
Before the shift key, typewriters had to have a separate key and typebar for upper-case letters; in essence, the typewriter had two keyboards, one above the other. With the shift key, manufacturing costs (and therefore purchase price) were greatly reduced, and typist operation was simplified; both factors contributed greatly to mass adoption of the technology.
Three-bank typewriters
Certain models further reduced the number of keys and typebars by making each key perform three functions – each typebar could type three different characters. These little three-row machines were portable and could be used by journalists.
Such three-row machines were popular with WWI journalists because they were lighter and more compact than four-bank typewriters, while they could type just as fast and use just as many symbols.
Such three-row machines, such as the Bar-Let and the Corona No. 3 Typewriter have two separate shift keys, a "CAP" shift (for uppercase) and a "FIG" shift (for numbers and symbols).
The Murray code was developed for a teletypewriter with a similar three-row typewriter keyboard.
Tab key
To facilitate typewriter use in business settings, a tab (tabulator) key was added in the late nineteenth century. Before using the key, the operator had to set mechanical "tab stops", pre-designated locations to which the carriage would advance when the tab key was pressed. This facilitated the typing of columns of numbers, freeing the operator from the need to manually position the carriage. The first models had one tab stop and one tab key; later ones allowed as many stops as desired, and sometimes had multiple tab keys, each of which moved the carriage a different number of spaces ahead of the decimal point (the tab stop), to facilitate the typing of columns with numbers of different length ($1.00, $10.00, $100.00, etc.)
Dead keys
Languages such as French, Spanish, and German required diacritics, special signs attached to or on top of the base letter: for example, a combination of the acute accent plus produced ; plus produced . In metal typesetting, , , and others were separate sorts. With mechanical typewriters, the number of whose characters (sorts) was constrained by the physical limits of the machine, the number of keys required was reduced by the use of dead keys. Diacritics such as (acute accent) would be assigned to a dead key, which did not move the platen forward, permitting another character to be imprinted at the same location; thus a single dead key such as the acute accent could be combined with ,,, and to produce ,,, and , reducing the number of sorts needed from 5 to 1. The typebars of "normal" characters struck a rod as they moved the metal character desired toward the ribbon and platen, and each rod depression moved the platen forward the width of one character. Dead keys had a typebar shaped so as not to strike the rod.
Character sizes
In English-speaking countries, ordinary typewriters printing fixed-width characters were standardized to print six horizontal lines per vertical inch, and had either of two variants of character width, one called pica for ten characters per horizontal inch and the other elite, for twelve. This differed from the use of these terms in printing, where pica is a linear unit (approximately of an inch) used for any measurement, the most common one being the height of a typeface.
Color
Some ribbons were inked in black and red stripes, each being half the width and running the entire length of the ribbon. A lever on most machines allowed switching between colors, which was useful for bookkeeping entries where negative amounts were highlighted in red. The red color was also used on some selected characters in running text, for emphasis. When a typewriter had this facility, it could still be fitted with a solid black ribbon; the lever was then used to switch to fresh ribbon when the first stripe ran out of ink. Some typewriters also had a third position which stopped the ribbon being struck at all. This enabled the keys to hit the paper unobstructed, and was used for cutting stencils for stencil duplicators (aka mimeograph machines).
"Noiseless" designs
The first typewriter to have the sliding type bars (laid out horizontally like a fan) that enable a typewriter to be "noiseless" was the American made Rapid which appeared briefly on the market in 1890. The Rapid also had the remarkable ability for the typist to have entire control of the carriage by manipulation of the keyboard alone. The two keys that achieve this are positioned at the top of the keyboard (seen in the detail image below). They are a "Lift" key that advances the paper, on the platen, to the next line and a "Return" key that causes the carriage to automatically swing back to the right, ready for one to type the new line. So an entire page could be typed without one's hands leaving the keyboard.
In the early part of the 20th century, a typewriter was marketed under the name Noiseless and advertised as "silent". It was developed by Wellington Parker Kidder and the first model was marketed by the Noiseless Typewriter Company in 1917. Noiseless portables sold well in the 1930s and 1940s, and noiseless standards continued to be manufactured until the 1960s.
In a conventional typewriter the type bar reaches the end of its travel simply by striking the ribbon and paper. The Noiseless, developed by Kidder, has a complex lever mechanism that decelerates the type bar mechanically before pressing it against the ribbon and paper in an attempt to dampen the noise.
Electric designs
Although electric typewriters would not achieve widespread popularity until nearly a century later, the basic groundwork for the electric typewriter was laid by the Universal Stock Ticker, invented by Thomas Edison in 1870. This device remotely printed letters and numbers on a stream of paper tape from input generated by a specially designed typewriter at the other end of a telegraph line.
Early electric models
Some electric typewriters were patented in the 19th century, but the first machine known to be produced in series is the Cahill of 1900.
Another electric typewriter was produced by the Blickensderfer Manufacturing Company, of Stamford, Connecticut, in 1902. Like the manual Blickensderfer typewriters, it used a cylindrical typewheel rather than individual typebars. The machine was produced in several variants but apparently not a commercial success, having come to market ahead of its time, before ubiquitous electrification.
The next step in the development of the electric typewriter came in 1910, when Charles and Howard Krum filed a patent for the first practical teletypewriter. The Krums' machine, named the Morkrum Printing Telegraph, used a typewheel rather than individual typebars. This machine was used for the first commercial teletypewriter system on Postal Telegraph Company lines between Boston and New York City in 1910.
James Fields Smathers of Kansas City invented what is considered the first practical power-operated typewriter in 1914. In 1920, after returning from Army service, he produced a successful model and in 1923 turned it over to the Northeast Electric Company of Rochester for development. Northeast was interested in finding new markets for their electric motors and developed Smathers's design so that it could be marketed to typewriter manufacturers, and from 1925 Remington Electric typewriters were produced powered by Northeast's motors.
After some 2,500 electric typewriters had been produced, Northeast asked Remington for a firm contract for the next batch. However, Remington was engaged in merger talks, which would eventually result in the creation of Remington Rand and no executives were willing to commit to a firm order. Northeast instead decided to enter the typewriter business for itself, and in 1929 produced the first Electromatic Typewriter.
In 1928, Delco, a division of General Motors, purchased Northeast Electric, and the typewriter business was spun off as Electromatic Typewriters, Inc. In 1933, Electromatic was acquired by IBM, which then spent $1 million on a redesign of the Electromatic Typewriter, launching the IBM Electric Typewriter Model 01.
In 1931, an electric typewriter was introduced by Varityper Corporation. It was called the Varityper, because a narrow cylinder-like wheel could be replaced to change the font.
In 1941, IBM announced the Electromatic Model 04 electric typewriter, featuring the revolutionary concept of proportional spacing. By assigning varied rather than uniform spacing to different sized characters, the Type 4 recreated the appearance of a typeset page, an effect that was further enhanced by including the 1937 innovation of carbon-film ribbons that produced clearer, sharper words on the page.
IBM Selectric
IBM introduced the IBM Selectric typewriter in 1961, which replaced the typebars with a spherical element (or typeball) slightly smaller than a golf ball, with reverse-image letters molded into its surface. The Selectric used a system of latches, metal tapes, and pulleys driven by an electric motor to rotate the ball into the correct position and then strike it against the ribbon and platen. The typeball moved laterally in front of the paper, instead of the previous designs using a platen-carrying carriage moving the paper across a stationary print position.
Due to the physical similarity, the typeball was sometimes referred to as a "golfball". The typeball design had many advantages, especially the elimination of "jams" (when more than one key was struck at once and the typebars became entangled) and in the ability to change the typeball, allowing multiple fonts to be used in a single document.
The IBM Selectric became a commercial success, dominating the office typewriter market for at least two decades. IBM also gained an advantage by marketing more heavily to schools than did Remington, with the idea that students who learned to type on a Selectric would later choose IBM typewriters over the competition in the workplace as businesses replaced their old manual models.
Later models of IBM Executives and Selectrics replaced inked fabric ribbons with "carbon film" ribbons that had a dry black or colored powder on a clear plastic tape. These could be used only once, but later models used a cartridge that was simple to replace. A side effect of this technology was that the text typed on the machine could be easily read from the used ribbon, raising issues where the machines were used for preparing classified documents (ribbons had to be accounted for to ensure that typists did not carry them from the facility).
A variation known as "Correcting Selectrics" introduced a correction feature, later imitated by competing machines, where a sticky tape in front of the carbon film ribbon could remove the black-powdered image of a typed character, eliminating the need for little bottles of white dab-on correction fluid and for hard erasers that could tear the paper. These machines also introduced selectable "pitch" so that the typewriter could be switched between pica type (10 characters per inch) and elite type (12 per inch), even within one document. Even so, all Selectrics were monospaced – each character and letterspace was allotted the same width on the page, from a capital "W" to a period. IBM did produce a successful typebar-based machine with five levels of proportional spacing, called the IBM Executive.
The only fully electromechanical Selectric Typewriter with fully proportional spacing and which used a Selectric type element was the expensive Selectric Composer, which was capable of right-margin justification (typing each line twice was required, once to calculate and again to print) and was considered a typesetting machine rather than a typewriter. Composer typeballs physically resembled those of the Selectric typewriter but were not interchangeable.In addition to its electronic successors, the Magnetic Tape Selectric Composer (MT/SC), the Mag Card Selectric Composer, and the Electronic Selectric Composer, IBM also made electronic typewriters with proportional spacing using the Selectric element that were considered typewriters or word processors instead of typesetting machines.
The first of these was the relatively obscure Mag Card Executive, which used 88-character elements. Later, some of the same typestyles used for it were used on the 96-character elements used on the IBM Electronic Typewriter 50 and the later models 65 and 85.
By 1970, as offset printing began to replace letterpress printing, the Composer would be adapted as the output unit for a phototypesetting system. The system included a computer-driven input station to capture the key strokes on magnetic tape and insert the operator's format commands, and a Composer unit to read the tape and produce the formatted text for photo reproduction.
The IBM 2741 terminal was a popular example of a Selectric-based computer terminal, and similar mechanisms were employed as the console devices for many IBM System/360 computers. These mechanisms used "ruggedized" designs compared to those in standard office typewriters.
Later electric models
Some of IBM's advances were later adopted in less expensive machines from competitors. For example, Smith-Corona electric typewriters introduced in 1973 switched to interchangeable Coronamatic (SCM-patented) ribbon cartridges. including fabric, film, erasing, and two-color versions. At about the same time, the advent of photocopying meant that carbon copies, correction fluid and erasers were less and less necessary; only the original need be typed, and photocopies made from it.
Electronic typewriters
The final major development of the typewriter was the electronic typewriter. Most of these replaced the typeball with a plastic or metal daisy wheel mechanism (a disk with the letters molded on the outside edge of the "petals"). The daisy wheel concept first emerged in printers developed by Diablo Systems in the 1970s. The first electronic daisywheel typewriter marketed in the world (in 1976) is the Olivetti Tes 501, and subsequently in 1978, the Olivetti ET101 (with function display) and Olivetti TES 401 (with text display and floppy disk for memory storage). This has allowed Olivetti to maintain the world record in the design of electronic typewriters, proposing increasingly advanced and performing models in the following years.
Unlike the Selectrics and earlier models, these really were "electronic" and relied on integrated circuits and electromechanical components. These typewriters were sometimes called display typewriters, dedicated word processors or word-processing typewriters, though the latter term was also frequently applied to less sophisticated machines that featured only a tiny, sometimes just single-row display. Sophisticated models were also called word processors, though today that term almost always denotes a type of software program. Manufacturers of such machines included Olivetti (TES501, first totally electronic Olivetti word processor with daisywheel and floppy disk in 1976; TES621 in 1979 etc.), Brother (Brother WP1 and WP500 etc., where WP stood for word processor), Canon (Canon Cat), Smith-Corona (PWP, i.e. Personal Word Processor line) and Philips/Magnavox (VideoWriter).
Decline
The pace of change was so rapid that it was common for clerical staff to have to learn several new systems, one after the other, in just a few years. While such rapid change is commonplace today, and is taken for granted, this was not always so; in fact, typewriting technology changed very little in its first 80 or 90 years.
Due to falling sales, IBM sold its typewriter division in 1991 to the newly formed Lexmark, completely exiting from a market it once dominated.
The increasing dominance of personal computers, desktop publishing, the introduction of low-cost, truly high-quality laser and inkjet printer technologies, and the pervasive use of web publishing, email, text messaging, and other electronic communication techniques have largely replaced typewriters in the United States. Still, , typewriters continued to be used by a number of government agencies and other institutions in the US, where they are primarily used to fill preprinted forms. According to a Boston typewriter repairman quoted by The Boston Globe, "Every maternity ward has a typewriter, as well as funeral homes".
A rather specialized market for typewriters exists due to the regulations of many correctional systems in the US, where prisoners are prohibited from having computers or telecommunication equipment, but are allowed to own typewriters. The Swintec corporation (headquartered in Moonachie, New Jersey), which, as of 2011, still produced typewriters at its overseas factories (in Japan, Indonesia, and/or Malaysia), manufactures a variety of typewriters for use in prisons, made of clear plastic (to make it harder for prisoners to hide prohibited items inside it). As of 2011, the company had contracts with prisons in 43 US states.
In April 2011, Godrej and Boyce, a Mumbai-based manufacturer of mechanical typewriters, closed its doors, leading to a flurry of news reports that the "world's last typewriter factory" had shut down. The reports were quickly contested, with opinions settling to agree that it was indeed the world's last producer of manual typewriters.
In November 2012, Brother's UK factory manufactured what it claimed to be the last typewriter ever made in the UK; the typewriter was donated to the London Science Museum.
Russian typewriters use Cyrillic, which has made the ongoing Azerbaijani reconversion from Cyrillic to Latin alphabet more difficult. In 1997, the government of Turkey offered to donate western typewriters to the Republic of Azerbaijan in exchange for more zealous and exclusive promotion of the Latin alphabet for the Azerbaijani language; this offer, however, was declined.
In Latin America and Africa, mechanical typewriters are still common because they can be used without electrical power. In Latin America, the typewriters used are most often Brazilian models; Brazil continues to produce mechanical (Facit) and electronic (Olivetti) typewriters to the present day.
The early 21st century saw revival of interest in typewriters among certain subcultures, including makers, steampunks, hipsters, and street poets.
Correction technologies
According to the standards taught in secretarial schools in the mid-20th century, a business letter was supposed to have no mistakes and no visible corrections.
Typewriter erasers
The traditional erasing method involved the use of a special typewriter eraser made of hard rubber that contained an abrasive material. Some were thin, flat disks, pink or gray, approximately in diameter by thick, with a brush attached from the center, while others looked like pink pencils, with a sharpenable eraser at the "lead" end and a stiff nylon brush at the other end. Either way, these tools made possible erasure of individual typed letters. Business letters were typed on heavyweight, high-rag-content bond paper, not merely to provide a luxurious appearance, but also to stand up to erasure.
Typewriter eraser brushes were necessary for clearing eraser crumbs and paper dust, and using the brush properly was an important element of typewriting skill; if erasure detritus fell into the typewriter, a small buildup could cause the typebars to jam in their narrow supporting grooves.
Erasing shield
Erasing a set of carbon copies was particularly difficult, and called for the use of a device called an erasing shield or eraser shield (a thin stainless-steel rectangle about with several tiny holes in it) to prevent the pressure of erasing on the upper copies from producing carbon smudges on the lower copies. To correct copies, typists had to go from one carbon copy layer to the next carbon copy layer, trying not to get their fingers dirty as they leafed through the carbon papers, and moving and repositioning the eraser shield and eraser for each copy.
Erasable bond
Paper companies produced a special form of typewriter paper called erasable bond (for example, Eaton's Corrasable Bond). This incorporated a thin layer of material that prevented ink from penetrating and was relatively soft and easy to remove from the page. An ordinary soft pencil eraser could quickly produce perfect erasures on this kind of paper. However, the same characteristics that made the paper erasable made the characters subject to smudging due to ordinary friction and deliberate alteration after the fact, making it unacceptable for business correspondence, contracts, or any archival use.
Correction fluid
In the 1950s and 1960s, correction fluid made its appearance, under brand names such as Liquid Paper, Wite-Out and Tipp-Ex; it was invented by Bette Nesmith Graham. Correction fluid was a kind of opaque, white, fast-drying paint that produced a fresh white surface onto which, when dry, a correction could be retyped. However, when held to the light, the covered-up characters were visible, as was the patch of dry correction fluid (which was never perfectly flat, and frequently not a perfect match for the color, texture, and luster of the surrounding paper). The standard trick for solving this problem was photocopying the corrected page, but this was possible only with high quality photocopiers.
A different fluid was available for correcting stencils. It sealed up the stencil ready for retyping but did not attempt to color match.
Legacy
Keyboard layouts
QWERTY
The 1874 Sholes & Glidden typewriters established the "QWERTY" layout for the letter keys. During the period in which Sholes and his colleagues were experimenting with this invention, other keyboard arrangements were apparently tried, but these are poorly documented. The QWERTY layout of keys has become the de facto standard for English-language typewriter and computer keyboards. Other languages written in the Latin alphabet sometimes use variants of the QWERTY layouts, such as the French AZERTY, the Italian QZERTY and the German QWERTZ layouts.
The QWERTY layout is not the most efficient layout possible for the English language. Touch-typists are required to move their fingers between rows to type the most common letters. Although the QWERTY keyboard was the most commonly used layout in typewriters, a better, less strenuous keyboard was being searched for throughout the late 1900s.
One popular but incorrect explanation for the QWERTY arrangement is that it was designed to reduce the likelihood of internal clashing of typebars by placing commonly used combinations of letters farther from each other inside the machine.
Other layouts for English
A number of radically different layouts such as Dvorak have been proposed to reduce the perceived inefficiencies of QWERTY, but none have been able to displace the QWERTY layout; their proponents claim considerable advantages, but so far none has been widely used. The Blickensderfer typewriter with its DHIATENSOR layout may have possibly been the first attempt at optimizing the keyboard layout for efficiency advantages.
On modern keyboards, the exclamation point is the shifted character on the 1 key, because these were the last characters to become "standard" on keyboards. Holding the spacebar down usually suspended the carriage advance mechanism (a so-called "dead key" feature), allowing one to superimpose multiple keystrikes on a single location. The ¢ symbol (meaning cents) was located above the number 6 on American electric typewriters, whereas ANSI-INCITS-standard computer keyboards have ^ instead.
Keyboards for other languages
The keyboards for other Latin languages are broadly similar to QWERTY but are optimised for the relevant orthography. In addition to some changes in the order of letters, perhaps the most obvious is the presence of precomposed characters and diacritics.
Many non-Latin alphabets have keyboard layouts that have nothing to do with QWERTY. The Russian layout, for instance, puts the common trigrams ыва, про, and ить on adjacent keys so that they can be typed by rolling the fingers.
Typewriters were also made for East Asian languages with thousands of characters, such as Chinese or Japanese. They were not easy to operate, but professional typists used them for a long time until the development of electronic word processors and laser printers in the 1980s.
Typewriter conventions
A number of typographical conventions stem from the typewriter's characteristics and limitations. For example, the QWERTY keyboard typewriter did not include keys for the en dash and the em dash. To overcome this limitation, users typically typed more than one adjacent hyphen to approximate these symbols. This typewriter convention is still sometimes used today, even though modern computer word processing applications can input the correct en and em dashes for each font type.
Other examples of typewriter practices that are sometimes still used in desktop publishing systems include inserting a double space between sentences, and the use of the typewriter apostrophe, , and straight quotes, , as quotation marks and prime marks. The practice of underlining text in place of italics and the use of all capitals to provide emphasis are additional examples of typographical conventions that derived from the limitations of the typewriter keyboard that still carry on today.
Many older typewriters did not include a separate key for the numeral or the exclamation point , and some even older ones also lacked the numeral zero, . Typists who trained on these machines learned the habit of using the lowercase letter ("ell") for the digit , and the uppercase ("oh") for the zero. A cents symbol, was created by combining (over-striking) a lower case with a slash character (typing , then backspace, then ). Similarly, the exclamation point was created by combining an apostrophe and a period ( ≈).
Terminology repurposed for the computer age
Some terminology from the typewriter age has survived into the computer era.
backspace (BS) – a keystroke that moved the cursor backwards one position (on a typewriter, this moved the physical platen backwards), to enable a character to be overtyped. Originally this was used to combine characters (for example, the sequence , backspace, to make ). Subsequently it facilitated "erase and retype" corrections (using correction tape or fluid.) Only the latter concept has survived into the computer age.
carriage return (CR) – return to the first column of text. (Most typewriters switched automatically to the next line. In computer systems, "line feed" (see below) is a function that is controlled independently.)
cursor – a marker used to indicate where the next character will be printed. The cursor was originally a term to describe the clear slider on a slide rule; on typewriters, it was the paper that moved and the insertion point was fixed.
cut and paste – taking text, a numerical table, or an image and pasting it into a document. The term originated when such compound documents were created using manual paste up techniques for typographic page layout. Actual brushes and paste were later replaced by hot-wax machines equipped with cylinders that applied melted adhesive wax to developed prints of "typeset" copy. This copy was then cut out with knives and rulers, and slid into position on layout sheets on slanting layout tables. After the "copy" had been correctly positioned and squared up using a T-square and set square, it was pressed down with a brayer, or roller. The whole point of the exercise was to create so-called "camera-ready copy" which existed only to be photographed and then printed, usually by offset lithography.
dead key – a key that, when typed, does not advance the typing position, thus allowing another character to be overstruck on top of the original character. This was typically used to combine diacritical marks with letters they modified (e.g. è can be generated by first pressing and then ). In Europe, where most languages have diacritics, a typical mechanical arrangement meant that hitting the accent key typed the symbol but did not advance the carriage, consequently the next character to be typed 'landed' on the same position. It was this method that carried across to the computer age whereas an alternative method (press the space bar simultaneously) did not.
line feed (LF), also called "newline" – Whereas most typewriters rolled the paper forward automatically on a "carriage return), this is an explicit control character on computer systems that moves the cursor to the next on-screen line of text. (But not to the beginning of that line a CR is also needed if that effect is desired.)
shift – a modifier key used to type capital letters and other alternate "upper case" characters; when pressed and held down, would shift a typewriter's mechanism to allow a different typebar impression (such as 'D' instead of 'd') to press into the ribbon and print on a page. The concept of a shift key or modifier key was later extended to Ctrl, Alt, AltGr and Super ("Windows" or "Apple") keys on modern computer keyboards. The generalized concept of a shift key reached its apex in the MIT space-cadet keyboard.
tab (HT), shortened from "horizontal tab" or "tabulator stop" – caused the print position to advance horizontally to the next pre-set "tab stop". This was used for typing lists and tables with vertical columns of numbers or words.
The vertical tab (VT) control character, named by analogy with HT, was designed for use with early computer line printers, and would cause the fan-fold paper to be fed until the next line's position.
tty, short for teletypewriter – used in Unix-like operating systems to designate a given "terminal".
Social effects
When Remington started marketing typewriters, the company assumed the machine would not be used for composing but for transcribing dictation, and that the person typing would be a woman. The 1800s Sholes and Glidden typewriter had floral ornamentation on the case.
During World Wars I and II, increasing numbers of women were entering the workforce. In the United States, women often started in the professional workplace as copy typists. Being a typist was considered the right choice for a "good girl", meaning women who present themselves as being chaste and having good conduct. According to the 1900 census, 94.9% of stenographers and typists were unmarried women. This also led to an increase in schools and classes for typing in order to prepare for a future job. Moreover, the word "typewriter" also became associated with the women who typed during the timeperiod.
Questions about morals made a salacious businessman making sexual advances to a female typist into a cliché of office life, appearing in vaudeville and movies. The "Tijuana bibles" – adult comic books produced in Mexico for the American market, starting in the 1930s – often featured women typists. In one panel, a businessman in a three-piece suit, ogling his secretary's thigh, says, "Miss Higby, are you ready for—ahem!—er—dictation?"
The typewriter was a useful machine during the censorship era of the Soviet government, starting during the Russian Civil War (1917–1922). Samizdat was a form of surreptitious self-publication used when the government was censoring what literature the public could see. The Soviet government signed a Decree on Press which prohibited the publishing of any written work that had not been previously officially reviewed and approved. Unapproved work was copied manually, most often on typewriters. In 1983, a new law required anyone who needed a typewriter to get police permission to buy or keep one. In addition, the owner would have to register a typed sample of all its letters and numbers, to ensure that any illegal literature typed with it could be traced back to its source. The typewriter became increasingly popular as the interest in prohibited books grew.
Writers with notable associations with typewriters
Early adopters
Henry James dictated to a typist.
Mark Twain claimed in his autobiography that he was the first important writer to present a publisher with a typewritten manuscript, for The Adventures of Tom Sawyer (1876). Research showed that Twain's memory was incorrect and that the first book submitted in typed form was Life on the Mississippi (1883, also by Twain).
Others
William S. Burroughs wrote in some of his novels – and possibly believed – that "a machine he called the 'Soft Typewriter' was writing our lives, and our books, into existence", according to a book review in The New Yorker. In the 1991 film adaptation of his 1959 novel Naked Lunch, his typewriter is a living, insect-like entity (voiced by North American actor Peter Boretski) and actually dictates the book to him.
J. R. R. Tolkien was accustomed to typing from awkward positions: "balancing his typewriter on his attic bed, because there was no room on his desk".
Jack Kerouac, a fast typist at 100 words per minute, typed his 1957 novel On the Road on a roll of paper so he would not be interrupted by having to change the paper. Within two weeks of starting to write On the Road, Kerouac had one single-spaced paragraph, long. Some scholars say the scroll was shelf paper; others contend it was a Thermal-fax roll; another theory is that the roll consisted of sheets of architect's paper taped together. Kerouac himself stated that he used rolls of teletype paper.
Don Marquis purposely used the limitations of a typewriter (or more precisely, a particular typist) in his archy and mehitabel series of newspaper columns, which were later compiled into a series of books. According to his literary conceit, a cockroach named "Archy" was a reincarnated free-verse poet, who would type articles overnight by jumping onto the keys of a manual typewriter. The writings were typed completely in lower case, because of the cockroach's inability to generate the heavy force needed to operate the shift key. The lone exception is the poem "CAPITALS AT LAST" from archys life of mehitabel, written in 1933.
Late users
Richard Polt, a philosophy professor at Xavier University in Cincinnati who collects typewriters, has edited ETCetera, a quarterly magazine about historic writing machines, and is the author of the book The Typewriter Revolution: A Typist's Companion for the 21st Century.
William Gibson used a Hermes 2000 model manual typewriter to write his 1984 novel Neuromancer and half of Count Zero (1983) before a mechanical failure and lack of replacement parts forced him to upgrade to an Apple IIc computer.
Harlan Ellison used typewriters for his entire career, and when he was no longer able to have them repaired, learned to do it himself; he repeatedly stated his belief that computers are bad for writing, maintaining that "Art is not supposed to be easier!"
Cormac McCarthy wrote his novels on an Olivetti Lettera 32 typewriter until his death. In 2009, the Lettera he obtained from a pawn shop in 1963, on which nearly all his novels and screenplays were written, was auctioned for charity at Christie's for US$254,500; McCarthy obtained an identical replacement for $20 to continue writing on.
Will Self explains why he uses a manual typewriter: "I think the computer user does their thinking on the screen, and the non-computer user is compelled, because he or she has to retype a whole text, to do a lot more thinking in the head."
Ted Kaczynski (the "Unabomber") infamously used two old manual typewriters to write his polemic essays and messages.
Actor Tom Hanks uses and collects manual typewriters. To control the size of his collection, he gifts autographed machines to appreciative fans and repair shops around the world.
Historian David McCullough used a Royal typewriter to compose his books.
Biographer Robert Caro has used various models of the Smith Corona Electra 210 to write his biographies of Robert Moses and Lyndon Johnson.
Typewriters in popular culture
In music
Erik Satie's 1917 score for the ballet Parade includes a "Mach. à écrire" as a percussion instrument, along with (elsewhere) a roulette wheel and a pistol.
The composer Leroy Anderson wrote The Typewriter (1950) for orchestra and typewriter, and it has since been used as the theme for numerous radio programs. The solo instrument is a real typewriter played by a percussionist. The piece was later made famous by comedian Jerry Lewis as part of his regular routine both on screen and stage, most notably in the 1963 film Who's Minding the Store?.
A typewriter plays an integral part (and is used on stage as a prop) in the song 'Opening Doors', from Stephen Sondheim's musical Merrily We Roll Along (1981).
Wordy Rappinghood, a 1981 single by Tom Tom Club, opens with the sound of a typewriter.
Typewriter samples are woven into the texture of 'Dissidents', the opening track of Thomas Dolby's 1984 album The Flat Earth.
The Boston Typewriter Orchestra (BTO), a comedic musical percussion group, has performed at numerous art festivals, clubs, and parties since 2004.
Max Richter's The Blue Notebooks (2004) features the sound of the typewriter underneath the narration of Tilda Swinton.
South Korean improviser Ryu Hankil frequently performs on typewriters, most prominently in his 2009 album Becoming Typewriter.
Other
The 2012 French comedy movie Populaire, starring Romain Duris and Déborah François, centers on a young secretary in the 1950s striving to win typewriting speed competitions.
The manga (2015–2020) and anime (2018) Violet Evergarden series follows a disabled war veteran who learns to type because her handwriting has been impaired, and soon she becomes a popular typist.
California Typewriter, a 2016 documentary film, investigates the culture of typewriter enthusiasts, including an eponymous repair store in Berkeley, California.
Forensic examination
Typewritten documents may be examined by forensic document examiners. This is done primarily to determine 1) the make and/or model of the typewriter used to produce a document, or 2) whether or not a particular suspect typewriter might have been used to produce a document.
The determination of a make and/or model of typewriter is a 'classification' problem and several systems have been developed for this purpose. These include the original Haas Typewriter Atlases (Pica version) and (Non-Pica version) and the TYPE system developed by Philip Bouffard, the Royal Canadian Mounted Police's Termatrex Typewriter classification system, and Interpol's typewriter classification system, among others.
The earliest reference in fictional literature to the potential identification of a typewriter as having produced a document was by Sir Arthur Conan Doyle, who wrote the Sherlock Holmes short story "A Case of Identity" in 1891.
In non-fiction, the first document examiner to describe how a typewriter might be identified was William E. Hagan who wrote, in 1894, "All typewriter machines, even when using the same kind of type, become more or less peculiar by use as to the work done by them". Other early discussions of the topic were provided by A. S. Osborn in his 1908 treatise, Typewriting as Evidence, and again in his 1929 textbook, Questioned Documents.
A modern description of the examination procedure is laid out in ASTM Standard E2494-08 (Standard Guide for Examination of Typewritten Items).
Typewriter examination was used in the Leopold and Loeb and Alger Hiss cases.
In the Eastern Bloc, typewriters (together with printing presses, copy machines, and later computer printers) were a controlled technology, with secret police in charge of maintaining records of the typewriters and their owners. In the Soviet Union, the First Department of each organization sent data on organization's typewriters to the KGB. This posed a significant risk for dissidents and samizdat authors. In Romania, according to State Council Decree No. 98 of March 28, 1983, owning a typewriter, both by businesses or by private persons, was subject to an approval given by the local police authorities. People previously convicted of any crime or those who because of their behaviour were considered to be "a danger to public order or to the security of the state" were refused approval. In addition, once a year, typewriter owners had to take the typewriter to the local police station, where they would be asked to type a sample of all the typewriter's characters. It was also forbidden to borrow, lend, or repair typewriters other than at the places that had been authorized by the police.
Collections
Public and private collections of typewriters exist around the world, including:
Schreibmaschinenmuseum Peter Mitterhofer (Parcines, Italy)
Museo della Macchina da Scrivere (Milan, Italy)
Liverpool Typewriter Museum (Liverpool, England)
Museum of Printing – MoP (Haverhill, Massachusetts, US)
Chestnut Ridge Typewriter Museum (Fairmont, West Virginia, US)
Technical Museum of the Empordà (Figueres, Girona, Spain)
Musée de la machine à écrire (Lausanne, Switzerland)
Lu Hanbin Typewriter Museum Shanghai (Shanghai, China)
Wattens Typewriter Museum (Wattens, Austria)
German Typewriter Museum (Bayreuth, Germany)
Tayfun Talipoğlu Typewriter Museum (Odunpazarı, Eskişehir, Turkey)
Several online-only virtual museums collect and display information about typewriters and their history:
Virtual Typewriter Museum
Chuck & Rich's Antique Typewriter Website
Mr. Martin's Typewriter Museum
Gallery
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44371 | https://en.wikipedia.org/wiki/Triton%20%28moon%29 | Triton (moon) | Triton is the largest natural satellite of the planet Neptune. It is the only moon of Neptune massive enough to be rounded under its own gravity and hosts a thin, hazy atmosphere. Triton orbits Neptune in a retrograde orbit—revolving in the opposite direction to the parent planet's rotation—the only large moon in the Solar System to do so. Triton is thought to have once been a dwarf planet from the Kuiper belt, captured into Neptune's orbit by the latter's gravity.
At in diameter, Triton is the seventh-largest moon in the Solar System, the second-largest planetary moon in relation to its primary (after Earth's Moon), and larger than all of the known dwarf planets. The mean density is , reflecting a composition of approximately 30–45% water ice by mass, with the rest being mostly rock and metal. Triton is differentiated, with a crust of primarily ice atop a probable subsurface ocean of liquid water and a solid rocky-metallic core at its center. Although Triton's orbit is nearly circular with a very low orbital eccentricity of , Triton's interior may still experience tidal heating through obliquity tides.
Triton is one of the most geologically active worlds in the Solar System, with an estimated average surface age of less than 100 million years old. Its surface is covered by frozen nitrogen and is geologically young, with very few impact craters. Young, intricate cryovolcanic and tectonic terrains suggest a complex geological history. The atmosphere of Triton is composed primarily of nitrogen, with minor components of methane and carbon monoxide. Triton's atmosphere is relatively thin and strongly variable, with its atmospheric surface pressure varying by up to a factor of three within the past 30 years. Triton's atmosphere supports clouds of nitrogen ice crystals and a layer of organic atmospheric haze.
Triton was the first Neptunian moon to be discovered, on October 10, 1846, by English astronomer William Lassell. The 1989 flyby of Triton by the Voyager 2 spacecraft remains the only up-close visit to the moon as of 2025. As the probe was only able to study about 40% of the moon's surface, multiple concept missions have been developed to revisit Triton. These include a Discovery-class Trident and New Frontiers-class Triton Ocean Worlds Surveyor and Nautilus.
Discovery and naming
Triton was discovered by British astronomer William Lassell on October 10, 1846, just 17 days after the discovery of Neptune. When John Herschel received news of Neptune's discovery, he wrote to Lassell suggesting he search for possible moons. Lassell discovered Triton eight days later. Lassell also claimed for a period to have discovered rings. Although Neptune was later confirmed to have rings, they are so faint and dark that it is not plausible he saw them. A brewer by trade, Lassell spotted Triton with his self-built aperture metal mirror reflecting telescope (also known as the "two-foot" reflector). This telescope was donated to the Royal Observatory, Greenwich in the 1880s, but was eventually dismantled.
Triton is named after the Greek sea god Triton (Τρίτων), the son of Poseidon (the Greek god corresponding to the Roman Neptune). The name was first proposed by Camille Flammarion in his 1880 book Astronomie Populaire, and was officially adopted many decades later. Until the discovery of the second moon Nereid in 1949, Triton was commonly referred to as "the satellite of Neptune". Lassell did not name his discovery; he later successfully suggested the name Hyperion for the eighth moon of Saturn when he discovered it.
Orbit and rotation
Triton is unique among all large moons in the Solar System for its retrograde orbit around its planet (i.e. it orbits in a direction opposite to the planet's rotation). Most of the outer irregular moons of Jupiter and Saturn also have retrograde orbits, as do some of the irregular moons of Uranus and Neptune. However, these moons are all much more distant from their primaries and are small in comparison with the largest of them (Phoebe) having only 8% of the diameter (and 0.03% of the mass) of Triton.
Triton's orbit is associated with two tilts, the obliquity of Neptune's rotation to Neptune's orbit, 30°, and the inclination of Triton's orbit to Neptune's rotation, 157° (an inclination over 90° indicates retrograde motion). Triton's orbit precesses forward relative to Neptune's rotation with a period of about 678 Earth years (4.1 Neptunian years), making its Neptune-orbit-relative inclination vary between 127° and 173°. That inclination is currently 130°; Triton's orbit is now near its maximum departure from coplanarity with Neptune's.
Triton's rotation is tidally locked to be synchronous with its orbit around Neptune: it keeps one face oriented toward the planet at all times. Its equator is almost exactly aligned with its orbital plane. At present, Triton's rotational axis is about 40° from Neptune's orbital plane, hence as Neptune orbits the Sun, Triton's polar regions take turns facing the Sun, resulting in seasonal changes as one pole, then the other moves into the sunlight. Such changes were observed in 2010.
Triton's revolution around Neptune has become a nearly perfect circle with an eccentricity of almost zero. Viscoelastic damping from tides alone is not thought to be capable of circularizing Triton's orbit in the time since the origin of the system, and gas drag from a prograde debris disc is likely to have played a substantial role. Tidal interactions also cause Triton's orbit, which is already closer to Neptune than the Moon is to Earth, to gradually decay further; predictions are that 3.6 billion years from now, Triton will pass within Neptune's Roche limit. This will result in either a collision with Neptune's atmosphere or the breakup of Triton, forming a new ring system similar to that found around Saturn.
Capture
The current understanding of moons in retrograde orbits means they cannot form in the same region of the solar nebula as the planets they orbit. Therefore Triton must have been captured from elsewhere in the Solar System. Astrophysicists believe it might have originated in the Kuiper belt, a ring of small icy objects extending from just inside the orbit of Neptune to about 50 AU from the Sun. Thought to be the point of origin for the majority of short-period comets observed from Earth, the belt is also home to several large, planet-like bodies including Pluto, which is now recognized as the largest in a population of Kuiper belt objects (the plutinos) locked in resonant orbits with Neptune. Triton is only slightly larger than Pluto and is nearly identical in composition, which has led to the hypothesis that the two share a common origin.
This has been further supported in a 2024 study of the chemical composition of Pluto and Triton which suggests they originated in the same region of the outer Solar System before the latter was pulled into Neptune’s orbit. Kathleen Mandt at NASA's Goddard Space Flight Center in Maryland and her colleagues hypothesize that Triton and Pluto formed close to each other before the Solar System settled down. "They probably formed in the same region, which wouldn't be where the Kuiper belt is now—it would have been either closer or further away," says Mandt.
Studying prior data on the two bodies, the team found that both have a large amount of nitrogen and trace amounts of methane and carbon monoxide, which could have accumulated in the outer regions of the young nebula "For some reason, Triton was then ejected from this region and ensnared by Neptune". "They had to have formed beyond the water-ice line," says Mandt, referring to the distance from the sun where water would freeze into ice or snow, which is why Triton and Pluto have similar amounts of certain key elements. "One possibility is that the giant planets moved closer to the sun early in the first 100 million years or so of the Solar System, which may have disrupted the orbits of some bodies like Triton.", says Mandt.
The proposed capture of Triton may explain several features of the Neptunian system, including the extremely eccentric orbit of Neptune's moon Nereid and the scarcity of moons as compared to the other giant planets. Triton's initially eccentric orbit would have intersected the orbits of irregular moons and disrupted those of smaller regular moons, dispersing them through gravitational interactions.
Triton's eccentric post-capture orbit would have also resulted in tidal heating of its interior, which could have kept Triton fluid for a billion years; this inference is supported by evidence of differentiation in Triton's interior. This source of internal heat disappeared following tidal locking and circularization of the orbit.
Two types of mechanisms have been proposed for Triton's capture. To be gravitationally captured by a planet, a passing body must lose sufficient energy to be slowed down to a speed less than that required to escape. An early model of how Triton may have been slowed was by collision with another object, either one that happened to be passing by Neptune (which is unlikely), or a moon or proto-moon in orbit around Neptune (which is more likely). A more recent hypothesis suggests that, before its capture, Triton was part of a binary system. When this binary encountered Neptune, it interacted in such a way that the binary dissociated, with one portion of the binary expelled, and the other, Triton, becoming bound to Neptune. This event is more likely for more massive companions. This hypothesis is supported by several lines of evidence, including binaries being very common among the large Kuiper belt objects. The event was brief but gentle, saving Triton from collisional disruption. Events like this may have been common during the formation of Neptune, or later when it migrated outward.
However, simulations in 2017 showed that after Triton's capture, and before its orbital eccentricity decreased, it probably did collide with at least one other moon, and caused collisions between other moons.
Physical characteristics
Triton is the seventh-largest moon and sixteenth-largest object in the Solar System and is modestly larger than the dwarf planets Pluto and Eris. It is also the largest retrograde moon in the Solar System. It accounts for more than 99.5% of all the mass known to orbit Neptune, including the planet's rings and fifteen other known moons, and is also more massive than all known moons in the Solar System smaller than itself combined. Also, with a diameter 5.5% that of Neptune, it is the largest moon of a gas giant relative to its planet in terms of diameter, although Titan is bigger relative to Saturn in terms of mass (the ratio of Triton's mass to that of Neptune is approximately 1:4788). It has a radius, density (2.061 g/cm3), temperature and chemical composition similar to that of Pluto.
Triton's surface is covered with a transparent layer of annealed frozen nitrogen. Only 40% of Triton's surface has been observed and studied, but it may be entirely covered in such a thin sheet of nitrogen ice. Triton's surface consists of 55% nitrogen ice with other ices mixed in. Water ice comprises 15–35% and frozen carbon dioxide (dry ice) the remaining 10–20%. Trace ices include 0.1% methane and 0.05% carbon monoxide. There could also be ammonia ice on the surface, as there are indications of ammonia dihydrate in the lithosphere. Triton's mean density implies that it probably consists of about 30–45% water ice (including relatively small amounts of volatile ices), with the remainder being rocky material. Triton's surface area is 23 million km2, which is 4.5% of Earth, or 15.5% of Earth's land area. Triton has an unusually high albedo, reflecting 60–95% of the sunlight that reaches it, and it has changed only slightly since the first observations. By comparison, the Moon reflects only 11%. This high albedo causes Triton to reflect a lot of whatever little sunlight there is instead of absorbing it, causing it to have the coldest recorded temperature in the Solar System at . Triton's reddish color is thought to be the result of methane ice, which is converted to tholins under exposure to ultraviolet radiation.
Because Triton's surface indicates a long history of melting, models of its interior posit that Triton is differentiated, like Earth, into a solid core, a mantle and a crust. Water, the most abundant volatile in the Solar System, comprises Triton's mantle, enveloping a core of rock and metal. There is enough rock in Triton's interior for radioactive decay to maintain a liquid subsurface ocean to this day, similar to what is thought to exist beneath the surface of Europa and several other icy outer Solar System worlds. This is not thought to be adequate to power convection in Triton's icy crust. However, the strong obliquity tides are believed to generate enough additional heat to accomplish this and produce the observed signs of recent surface geological activity. The black material ejected is suspected to contain organic compounds, and if liquid water is present on Triton, it has been speculated that this could make it habitable for some form of life.
Atmosphere
Triton has a tenuous but well-structured and global nitrogen atmosphere, with trace amounts of carbon monoxide and small amounts of methane near its surface. Like Pluto's atmosphere, the atmosphere of Triton is thought to result from the evaporation of nitrogen from its surface. Its surface temperature is at least because Triton's nitrogen ice is in the warmer, hexagonal crystalline state, and the phase transition between hexagonal and cubic nitrogen ice occurs at that temperature. An upper limit in the low 40s (K) can be set from vapor pressure equilibrium with nitrogen gas in Triton's atmosphere. This is colder than Pluto's average equilibrium temperature of . Triton's surface atmospheric pressure is only about .
Turbulence at Triton's surface creates a troposphere (a "weather region") rising to an altitude of 8 km. Streaks on Triton's surface left by geyser plumes suggest that the troposphere is driven by seasonal winds capable of moving material over a micrometer in size. Unlike other atmospheres, Triton's lacks a stratosphere and instead has a thermosphere from altitudes of 8 to 950 km and an exosphere above that. The temperature of Triton's upper atmosphere, at , is higher than that at its surface, due to heat absorbed from solar radiation and Neptune's magnetosphere. A haze permeates most of Triton's troposphere, thought to be composed largely of hydrocarbons and nitriles created by the action of sunlight on methane. Triton's atmosphere also has clouds of condensed nitrogen that lie between 1 and 3 km from its surface.
In 1997, observations from Earth were made of Triton's limb as it passed in front of stars. These observations indicated the presence of a denser atmosphere than was deduced from Voyager 2 data. Other observations have shown an increase in temperature by 5% from 1989 to 1998. These observations indicated Triton was approaching an unusually warm southern hemisphere summer season that happens only once every few hundred years. Hypotheses for this warming include a change of frost patterns on Triton's surface and a change in ice albedo, which would allow more heat to be absorbed. Another hypothesis argues that temperature changes are a result of the deposition of dark, red material from geological processes. Because Triton's Bond albedo is among the highest in the Solar System, it is sensitive to small variations in spectral albedo. Based on the increase in atmospheric pressure between 1989 and 1997, it is estimated that by 2010 Triton's atmospheric pressure may have increased to as much as 4 Pa. By 2017, however, Triton's atmospheric surface pressure had nearly returned to Voyager 2 levels; the cause for the rapid spike in atmospheric pressure between 1989 and 2017 remains unexplained.
Surface features
All detailed knowledge of the surface of Triton was acquired from a distance of 40,000 km by the Voyager 2 spacecraft during a single encounter in 1989. The 40% of Triton's surface imaged by Voyager 2 revealed blocky outcrops, ridges, troughs, furrows, hollows, plateaus, icy plains and a few craters. Triton is relatively flat; its observed topography never varies beyond a kilometer. The impact craters observed are concentrated almost entirely in Triton's leading hemisphere. Analysis of crater density and distribution has suggested that in geological terms, Triton's surface is extremely young, with regions varying from an estimated 50 million years old to just an estimated 6 million years old. Fifty-five percent of Triton's surface is covered with frozen nitrogen, with water ice comprising 15–35% and frozen CO2 forming the remaining 10–20%. The surface also has deposits of tholins, a dark, tarry slurry of various organic chemical compounds.
Cryovolcanism
One of the largest cryovolcanic features found on Triton is Leviathan Patera, a caldera-like feature roughly 100 km in diameter seen near the equator. Surrounding this caldera is a massive cryovolcanic plain, Cipango Planum, which is at least 490,000 km2 in area; assuming Leviathan Patera is the primary vent, Leviathan Patera is one of the largest volcanic or cryovolcanic constructs in the Solar System. This feature is also connected to two enormous cryolava lakes seen northwest of the caldera. Because the cryolava on Triton is believed to be primarily water ice with some ammonia, these lakes would qualify as stable bodies of surface liquid water while they were molten. This is the first place such bodies have been found apart from Earth, and Triton is the only icy body known to feature cryolava lakes, although similar cryomagmatic extrusions can be seen on Ariel, Ganymede, Charon, and Titan.
Plumes
The Voyager 2 probe in 1989 observed a handful of geyser-like eruptions of nitrogen gas or water and entrained dust from beneath the surface of Triton in plumes up to 8 km high. Triton is thus one of the few bodies in the Solar System on which active eruptions of some sort have been observed. The best-observed examples are the Hili plume and Mahilani plume (named after a Zulu water sprite and a Tongan sea spirit, respectively).
The precise mechanism behind Triton's plumes is debated; one hypothesis is that Triton's plumes are driven by solar heating underneath a transparent or translucent layer of nitrogen ice, creating a sort of "solid greenhouse effect". As solar radiation warms the darker material beneath, this causes a rapid increase in pressure as the nitrogen begins to sublimate until enough pressure accumulates for it to erupt through the translucent layer. This model is largely supported by the observation that Triton was near peak southern summer at the time of Voyager 2s flyby, ensuring its southern polar cap was receiving prolonged sunlight. If this were the case, CO2 geysers on Mars are thought to erupt from its south polar cap each spring in the same way.
However, the significant geological activity on Triton has led to alternative proposals that the plumes may be cryovolcanic in nature, rather than driven by solar radiation. A cryovolcanic origin better explains the estimated output of Triton’s plumes, which possibly exceeds . This is similar to that which is estimated for Enceladus's cryovolcanic plumes at . However, if Triton's plumes are cryovolcanically driven, it remains to be explained why they predominantly appear over its southern polar cap. Triton's high surface heat flux may directly melt or vaporize nitrogen ice at the base of its polar caps, creating 'hot spots' which break through the ice or move to the ice caps' margins, before erupting explosively.
Though only observed up close once by the Voyager 2 spacecraft, it is estimated that a plume eruption on Triton may last up to a year.
Polar cap, plains and ridges
Triton's south polar region is covered by a highly reflective cap of frozen nitrogen and methane sprinkled by impact craters and openings of geysers. Little is known about the north pole because it was on the night side during the Voyager 2 encounter, but it is thought that Triton must also have a north polar ice cap.
The high plains found on Triton's eastern hemisphere, such as Cipango Planum, cover over and blot out older features, and are therefore almost certainly the result of icy lava washing over the previous landscape. The plains are dotted with pits, such as Leviathan Patera, which are probably the vents from which this lava emerged. The composition of the lava is unknown, although a mixture of ammonia and water is suspected.
Four roughly circular "walled plains" have been identified on Triton. They are the flattest regions so far discovered, with a variance in altitude of less than 200 m. They are thought to have formed from the eruption of icy lava. The plains near Triton's eastern limb are dotted with black spots, the maculae. Some maculae are simple dark spots with diffuse boundaries, and others comprise a dark central patch surrounded by a white halo with sharp boundaries. The maculae typically have diameters of about 100 km and widths of the halos of between 20 and 30 km.
There are extensive ridges and valleys in complex patterns across Triton's surface, probably the result of freeze–thaw cycles. Many also appear to be tectonic and may result from an extension or strike-slip faulting. There are long double ridges of ice with central troughs bearing a strong resemblance to Europan lineae (although they have a larger scale), and which may have a similar origin, possibly shear heating from strike-slip motion along faults caused by diurnal tidal stresses experienced before Triton's orbit was fully circularized. These faults with parallel ridges expelled from the interior cross complex terrain with valleys in the equatorial region. The ridges and furrows, or sulci, such as Yasu Sulci, Ho Sulci, and Lo Sulci, are thought to be of intermediate age in Triton's geological history, and in many cases to have formed concurrently. They tend to be clustered in groups or "packets".
Cantaloupe terrain
Triton's western hemisphere consists of a strange series of fissures and depressions known as "cantaloupe terrain" because it resembles the skin of a cantaloupe melon. Although it has few craters, it is thought that this is the oldest terrain on Triton. It probably covers much of Triton's western half.
Cantaloupe terrain, which is mostly dirty water ice, is only known to exist on Triton. It contains depressions in diameter. The depressions (cavi) are probably not impact craters because they are all of the similar size and have smooth curves. The leading hypothesis for their formation is diapirism, the rising of "lumps" of less dense material through a stratum of denser material. Alternative hypotheses include formation by collapses, or by flooding caused by cryovolcanism.
Impact craters
Due to constant erasure and modification by ongoing geological activity, impact craters on Triton's surface are relatively rare. A census of Triton's craters imaged by Voyager 2 found only 179 that were incontestably of impact origin, compared with 835 observed for Uranus's moon Miranda, which has only three percent of Triton's surface area. The largest crater observed on Triton thought to have been created by an impact is a feature called Mazomba. Although larger craters have been observed, they are generally thought to be volcanic.
The few impact craters on Triton are almost all concentrated in the leading hemisphere—that facing the direction of the orbital motion—with the majority concentrated around the equator between 30° and 70° longitude, resulting from material swept up from orbit around Neptune. Because it orbits with one side permanently facing the planet, astronomers expect that Triton should have fewer impacts on its trailing hemisphere, due to impacts on the leading hemisphere being more frequent and more violent. Voyager 2 imaged only 40% of Triton's surface, so this remains uncertain. However, the observed cratering asymmetry exceeds what can be explained based on the impactor populations, and implies a younger surface age for the crater-free regions (≤ 6 million years old) than for the cratered regions (≤ 50 million years old).
Observation and exploration
The orbital properties of Triton were already determined with high accuracy in the 19th century. It was found to have a retrograde orbit, at a very high angle of inclination to the plane of Neptune's orbit. The first detailed observations of Triton were not made until 1930. Little was known about the satellite until Voyager 2 flew by in 1989.
Before the flyby of Voyager 2, astronomers suspected that Triton might have liquid nitrogen seas and a nitrogen/methane atmosphere with a density as much as 30% that of Earth. Like the famous overestimates of the atmospheric density of Mars, this proved incorrect. As with Mars, a denser atmosphere is postulated for its early history.
The first attempt to measure the diameter of Triton was made by Gerard Kuiper in 1954. He obtained a value of 3,800 km. Subsequent measurement attempts arrived at values ranging from 2,500 to 6,000 km, or from slightly smaller than the Moon (3,474.2 km) to nearly half the diameter of Earth. Data from the approach of Voyager 2 to Neptune on August 25, 1989, led to a more accurate estimate of Triton's diameter (2,706 km).
In the 1990s, various observations from Earth were made of the limb of Triton using the occultation of nearby stars, which indicated the presence of an atmosphere and an exotic surface. Observations in late 1997 suggest that Triton is heating up and the atmosphere has become significantly denser since Voyager 2 flew past in 1989.
New concepts for missions to the Neptune system to be conducted in the 2010s were proposed by NASA scientists on numerous occasions over the last decades. All of them identified Triton as being a prime target and a separate Triton lander comparable to the Huygens probe for Titan was frequently included in those plans. No efforts aimed at Neptune and Triton went beyond the proposal phase and NASA's funding for missions to the outer Solar System is currently focused on the Jupiter and Saturn systems. A proposed lander mission to Triton, called Triton Hopper, would mine nitrogen ice from the surface of Triton and process it to be used as a propellant for a small rocket, enabling it to fly or 'hop' across the surface. Another concept, involving a flyby, was formally proposed in 2019 as part of NASA's Discovery Program under the name Trident. Neptune Odyssey is a mission concept for a Neptune orbiter with a focus on Triton being studied beginning April 2021 as a possible large strategic science mission by NASA that would launch in 2033 and arrive at the Neptune system in 2049. Two lower-cost mission concepts were subsequently developed for the New Frontiers program: the first the following June and the second in 2023. The first is Triton Ocean World Surveyor, which would launch in 2031 and arrive in 2047, and the second is Nautilus, which would launch August 2042 and arrive in April 2057.
Maps
| Physical sciences | Solar System | null |
44401 | https://en.wikipedia.org/wiki/Brown%20dwarf | Brown dwarf | Brown dwarfs are substellar objects that have more mass than the biggest gas giant planets, but less than the least massive main-sequence stars. Their mass is approximately 13 to 80 times that of Jupiter ()not big enough to sustain nuclear fusion of ordinary hydrogen (1H) into helium in their cores, but massive enough to emit some light and heat from the fusion of deuterium (2H). The most massive ones (> ) can fuse lithium (7Li).
Astronomers classify self-luminous objects by spectral type, a distinction intimately tied to the surface temperature, and brown dwarfs occupy types M, L, T, and Y. As brown dwarfs do not undergo stable hydrogen fusion, they cool down over time, progressively passing through later spectral types as they age.
Their name comes not from the color of light they emit but from their falling between white dwarf stars and "dark" planets in size. To the naked eye, brown dwarfs would appear in different colors depending on their temperature. The warmest ones are possibly orange or red, while cooler brown dwarfs would likely appear magenta or black to the human eye. Brown dwarfs may be fully convective, with no layers or chemical differentiation by depth.
Though their existence was initially theorized in the 1960s, it was not until the mid-1990s that the first unambiguous brown dwarfs were discovered. As brown dwarfs have relatively low surface temperatures, they are not very bright at visible wavelengths, emitting most of their light in the infrared. However, with the advent of more capable infrared detecting devices, thousands of brown dwarfs have been identified. The nearest known brown dwarfs are located in the Luhman 16 system, a binary of L- and T-type brown dwarfs about from the Sun. Luhman 16 is the third closest system to the Sun after Alpha Centauri and Barnard's Star.
History
Early theorizing
The objects now called "brown dwarfs" were theorized by Shiv S. Kumar in the 1960s to exist and were originally called black dwarfs, a classification for dark substellar objects floating freely in space that were not massive enough to sustain hydrogen fusion. However, (a) the term black dwarf was already in use to refer to a cold white dwarf; (b) red dwarfs fuse hydrogen; and (c) these objects may be luminous at visible wavelengths early in their lives. Because of this, alternative names for these objects were proposed, including and substar. In 1975, Jill Tarter suggested the term "brown dwarf", using "brown" as an approximate color.
The term "black dwarf" still refers to a white dwarf that has cooled to the point that it no longer emits significant amounts of light. However, the time required for even the lowest-mass white dwarf to cool to this temperature is calculated to be longer than the current age of the universe; hence such objects are expected to not yet exist.
Early theories concerning the nature of the lowest-mass stars and the hydrogen-burning limit suggested that a population I object with a mass less than 0.07 solar masses () or a population II object less than would never go through normal stellar evolution and would become a completely degenerate star. The resulting brown dwarf star is sometimes called a failed star. The first self-consistent calculation of the hydrogen-burning minimum mass confirmed a value between 0.07 and 0.08 solar masses for population I objects.
Deuterium fusion
The discovery of deuterium burning down to () and the impact of dust formation in the cool outer atmospheres of brown dwarfs in the late 1980s brought these theories into question. However, such objects were hard to find because they emit almost no visible light. Their strongest emissions are in the infrared (IR) spectrum, and ground-based IR detectors were too imprecise at that time to readily identify any brown dwarfs.
Since then, numerous searches by various methods have sought these objects. These methods included multi-color imaging surveys around field stars, imaging surveys for faint companions of main-sequence dwarfs and white dwarfs, surveys of young star clusters, and radial velocity monitoring for close companions.
GD 165B and class L
For many years, efforts to discover brown dwarfs were fruitless. In 1988, however, a faint companion to the white dwarf star GD 165 was found in an infrared search of white dwarfs. The spectrum of the companion GD 165B was very red and enigmatic, showing none of the features expected of a low-mass red dwarf. It became clear that GD 165B would need to be classified as a much cooler object than the latest M dwarfs then known. GD 165B remained unique for almost a decade until the advent of the Two Micron All-Sky Survey (2MASS) in 1997, which discovered many objects with similar colors and spectral features.
Today, GD 165B is recognized as the prototype of a class of objects now called "L dwarfs".
Although the discovery of the coolest dwarf was highly significant at the time, it was debated whether GD 165B would be classified as a brown dwarf or simply a very-low-mass star, because observationally it is very difficult to distinguish between the two.
Soon after the discovery of GD 165B, other brown-dwarf candidates were reported. Most failed to live up to their candidacy, however, because the absence of lithium showed them to be stellar objects. True stars burn their lithium within a little over 100 Myr, whereas brown dwarfs (which can, confusingly, have temperatures and luminosities similar to true stars) will not. Hence, the detection of lithium in the atmosphere of an object older than 100 Myr ensures that it is a brown dwarf.
Gliese 229B and class T
The first class "T" brown dwarf was discovered in 1994 by Caltech astronomers Shrinivas Kulkarni, Tadashi Nakajima, Keith Matthews and Rebecca Oppenheimer, and Johns Hopkins scientists Samuel T. Durrance and David Golimowski. It was confirmed in 1995 as a substellar companion to Gliese 229. Gliese 229b is one of the first two instances of clear evidence for a brown dwarf, along with Teide 1. Confirmed in 1995, both were identified by the presence of the 670.8 nm lithium line. The latter was found to have a temperature and luminosity well below the stellar range.
Its near-infrared spectrum clearly exhibited a methane absorption band at 2 micrometres, a feature that had previously only been observed in the atmospheres of giant planets and that of Saturn's moon Titan. Methane absorption is not expected at any temperature of a main-sequence star. This discovery helped to establish yet another spectral class even cooler than L dwarfs, known as "T dwarfs", for which Gliese 229B is the prototype.
Teide 1 and class M
The first confirmed class "M" brown dwarf was discovered by Spanish astrophysicists Rafael Rebolo (head of the team), María Rosa Zapatero-Osorio, and Eduardo L. Martín in 1994. This object, found in the Pleiades open cluster, received the name Teide 1. The discovery article was submitted to Nature in May 1995, and published on 14 September 1995. Nature highlighted "Brown dwarfs discovered, official" on the front page of that issue.
Teide 1 was discovered in images collected by the IAC team on 6 January 1994 using the 80 cm telescope (IAC 80) at Teide Observatory, and its spectrum was first recorded in December 1994 using the 4.2 m William Herschel Telescope at Roque de los Muchachos Observatory (La Palma). The distance, chemical composition, and age of Teide 1 could be established because of its membership in the young Pleiades star cluster. Using the most advanced stellar and substellar evolution models at that moment, the team estimated for Teide 1 a mass of , which is below the stellar-mass limit. The object became a reference in subsequent young brown dwarf related works.
In theory, a brown dwarf below is unable to burn lithium by thermonuclear fusion at any time during its evolution. This fact is one of the lithium test principles used to judge the substellar nature of low-luminosity and low-surface-temperature astronomical bodies.
High-quality spectral data acquired by the Keck 1 telescope in November 1995 showed that Teide 1 still had the initial lithium abundance of the original molecular cloud from which Pleiades stars formed, proving the lack of thermonuclear fusion in its core. These observations fully confirmed that Teide 1 is a brown dwarf, as well as the efficiency of the spectroscopic lithium test.
For some time, Teide 1 was the smallest known object outside the Solar System that had been identified by direct observation. Since then, over 1,800 brown dwarfs have been identified, even some very close to Earth, like Epsilon Indi Ba and Bb, a pair of brown dwarfs gravitationally bound to a Sun-like star 12 light-years from the Sun, and Luhman 16, a binary system of brown dwarfs at 6.5 light-years from the Sun.
Theory
The standard mechanism for star birth is through the gravitational collapse of a cold interstellar cloud of gas and dust. As the cloud contracts, it heats due to the Kelvin–Helmholtz mechanism. Early in the process the contracting gas quickly radiates away much of the energy, allowing the collapse to continue. Eventually, the central region becomes sufficiently dense to trap radiation. Consequently, the central temperature and density of the collapsed cloud increase dramatically with time, slowing the contraction, until the conditions are hot and dense enough for thermonuclear reactions to occur in the core of the protostar. For a typical star, gas and radiation pressure generated by the thermonuclear fusion reactions within its core will support it against any further gravitational contraction. Hydrostatic equilibrium is reached, and the star will spend most of its lifetime fusing hydrogen into helium as a main-sequence star.
If, however, the initial mass of the protostar is less than about , normal hydrogen thermonuclear fusion reactions will not ignite in the core. Gravitational contraction does not heat the small protostar very effectively, and before the temperature in the core can increase enough to trigger fusion, the density reaches the point where electrons become closely packed enough to create quantum electron degeneracy pressure. According to the brown dwarf interior models, typical conditions in the core for density, temperature and pressure are expected to be the following:
This means that the protostar is not massive or dense enough ever to reach the conditions needed to sustain hydrogen fusion. The infalling matter is prevented, by electron degeneracy pressure, from reaching the densities and pressures needed.
Further gravitational contraction is prevented and the result is a brown dwarf that simply cools off by radiating away its internal thermal energy. Note that, in principle, it is possible for a brown dwarf to slowly accrete mass above the hydrogen burning limit without initiating hydrogen fusion. This could happen via mass transfer in a binary brown dwarf system.
High-mass brown dwarfs versus low-mass stars
Lithium is generally present in brown dwarfs and not in low-mass stars. Stars, which reach the high temperature necessary for fusing hydrogen, rapidly deplete their lithium. Fusion of lithium-7 and a proton occurs, producing two helium-4 nuclei. The temperature necessary for this reaction is just below that necessary for hydrogen fusion. Convection in low-mass stars ensures that lithium in the whole volume of the star is eventually depleted. Therefore, the presence of the lithium spectral line in a candidate brown dwarf is a strong indicator that it is indeed a substellar object.
Lithium test
The use of lithium to distinguish candidate brown dwarfs from low-mass stars is commonly referred to as the lithium test, and was pioneered by Rafael Rebolo, Eduardo Martín and Antonio Magazzu. However, lithium is also seen in very young stars, which have not yet had enough time to burn it all.
Heavier stars, like the Sun, can also retain lithium in their outer layers, which never get hot enough to fuse lithium, and whose convective layer does not mix with the core where the lithium would be rapidly depleted. Those larger stars are easily distinguishable from brown dwarfs by their size and luminosity.
Conversely, brown dwarfs at the high end of their mass range can be hot enough to deplete their lithium when they are young. Dwarfs of mass greater than can burn their lithium by the time they are half a billion years old; thus the lithium test is not perfect.
Atmospheric methane
Unlike stars, older brown dwarfs are sometimes cool enough that, over very long periods of time, their atmospheres can gather observable quantities of methane, which cannot form in hotter objects. Dwarfs confirmed in this fashion include Gliese 229B.
Iron, silicate and sulfide clouds
Main-sequence stars cool, but eventually reach a minimum bolometric luminosity that they can sustain through steady fusion. This luminosity varies from star to star, but is generally at least 0.01% that of the Sun. Brown dwarfs cool and darken steadily over their lifetimes; sufficiently old brown dwarfs will be too faint to be detectable.
Clouds are used to explain the weakening of the iron hydride (FeH) spectral line in late L-dwarfs. Iron clouds deplete FeH in the upper atmosphere, and the cloud layer blocks the view to lower layers still containing FeH. The later strengthening of this chemical compound at cooler temperatures of mid- to late T-dwarfs is explained by disturbed clouds that allows a telescope to look into the deeper layers of the atmosphere that still contains FeH. Young L/T-dwarfs (L2-T4) show high variability, which could be explained with clouds, hot spots, magnetically driven aurorae or thermochemical instabilities. The clouds of these brown dwarfs are explained as either iron clouds with varying thickness or a lower thick iron cloud layer and an upper silicate cloud layer. This upper silicate cloud layer can consist out of quartz, enstatite, corundum and/or fosterite. It is however not clear if silicate clouds are always necessary for young objects. Silicate absorption can be directly observed in the mid-infrared at 8 to 12 μm. Observations with Spitzer IRS have shown that silicate absorption is common, but not ubiquitous, for L2-L8 dwarfs. Additionally, MIRI has observed silicate absorption in the planetary-mass companion VHS 1256b.
Iron rain as part of atmospheric convection processes is possible only in brown dwarfs, and not in small stars. The spectroscopy research into iron rain is still ongoing, but not all brown dwarfs will always have this atmospheric anomaly. In 2013, a heterogeneous iron-containing atmosphere was imaged around the B component in the nearby Luhman 16 system.
For late T-type brown dwarfs only a few variable searches were carried out. Thin cloud layers are predicted to form in late T-dwarfs from chromium and potassium chloride, as well as several sulfides. These sulfides are manganese sulfide, sodium sulfide and zinc sulfide. The variable T7 dwarf 2M0050–3322 is explained to have a top layer of potassium chloride clouds, a mid layer of sodium sulfide clouds and a lower layer of manganese sulfide clouds. Patchy clouds of the top two cloud layers could explain why the methane and water vapor bands are variable.
At the lowest temperatures of the Y-dwarf WISE 0855-0714 patchy cloud layers of sulfide and water ice clouds could cover 50% of the surface.
Low-mass brown dwarfs versus high-mass planets
Like stars, brown dwarfs form independently, but, unlike stars, they lack sufficient mass to "ignite" hydrogen fusion. Like all stars, they can occur singly or in close proximity to other stars. Some orbit stars and can, like planets, have eccentric orbits.
Size and fuel-burning ambiguities
Brown dwarfs are all roughly the same radius as Jupiter. At the high end of their mass range (), the volume of a brown dwarf is governed primarily by electron-degeneracy pressure, as it is in white dwarfs; at the low end of the range (), their volume is governed primarily by Coulomb pressure, as it is in planets. The net result is that the radii of brown dwarfs vary by only 10–15% over the range of possible masses. Moreover, the mass–radius relationship shows no change from about one Saturn mass to the onset of hydrogen burning (), suggesting that from this perspective brown dwarfs are simply high-mass Jovian planets. This can make distinguishing them from planets difficult.
In addition, many brown dwarfs undergo no fusion; even those at the high end of the mass range (over ) cool quickly enough that after 10 million years they no longer undergo fusion.
Heat spectrum
X-ray and infrared spectra are telltale signs of brown dwarfs. Some emit X-rays; and all "warm" dwarfs continue to glow tellingly in the red and infrared spectra until they cool to planet-like temperatures (under ).
Gas giants have some of the characteristics of brown dwarfs. Like the Sun, Jupiter and Saturn are both made primarily of hydrogen and helium. Saturn is nearly as large as Jupiter, despite having only 30% the mass. Three of the giant planets in the Solar System (Jupiter, Saturn, and Neptune) emit much more (up to about twice) heat than they receive from the Sun. All four giant planets have their own "planetary" systems, in the form of extensive moon systems.
Current IAU standard
Currently, the International Astronomical Union considers an object above (the limiting mass for thermonuclear fusion of deuterium) to be a brown dwarf, whereas an object under that mass (and orbiting a star or stellar remnant) is considered a planet. The minimum mass required to trigger sustained hydrogen burning (about ) forms the upper limit of the definition.
It is also debated whether brown dwarfs would be better defined by their formation process rather than by theoretical mass limits based on nuclear fusion reactions. Under this interpretation brown dwarfs are those objects that represent the lowest-mass products of the star formation process, while planets are objects formed in an accretion disk surrounding a star. The coolest free-floating objects discovered, such as WISE 0855, as well as the lowest-mass young objects known, like PSO J318.5−22, are thought to have masses below , and as a result are sometimes referred to as planetary-mass objects due to the ambiguity of whether they should be regarded as rogue planets or brown dwarfs. There are planetary-mass objects known to orbit brown dwarfs, such as 2M1207b,2MASS J044144b and Oph 98 B.
The 13-Jupiter-mass cutoff is a rule of thumb rather than a quantity with precise physical significance. Larger objects will burn most of their deuterium and smaller ones will burn only a little, and the 13Jupiter-mass value is somewhere in between. The amount of deuterium burnt also depends to some extent on the composition of the object, specifically on the amount of helium and deuterium present and on the fraction of heavier elements, which determines the atmospheric opacity and thus the radiative cooling rate.
As of 2011 the Extrasolar Planets Encyclopaedia included objects up to 25 Jupiter masses, saying, "The fact that there is no special feature around in the observed mass spectrum reinforces the choice to forget this mass limit". As of 2016, this limit was increased to 60 Jupiter masses, based on a study of mass–density relationships.
The Exoplanet Data Explorer includes objects up to 24 Jupiter masses with the advisory: "The 13 Jupiter-mass distinction by the IAU Working Group is physically unmotivated for planets with rocky cores, and observationally problematic due to the sin i ambiguity." The NASA Exoplanet Archive includes objects with a mass (or minimum mass) equal to or less than 30 Jupiter masses.
Sub-brown dwarf
Objects below , called sub-brown dwarfs or planetary-mass brown dwarfs, form in the same manner as stars and brown dwarfs (i.e. through the collapse of a gas cloud) but have a mass below the limiting mass for thermonuclear fusion of deuterium.
Some researchers call them free-floating planets, whereas others call them planetary-mass brown dwarfs.
Role of other physical properties in the mass estimate
While spectroscopic features can help to distinguish between low-mass stars and brown dwarfs, it is often necessary to estimate the mass to come to a conclusion. The theory behind the mass estimate is that brown dwarfs with a similar mass form in a similar way and are hot when they form. Some have spectral types that are similar to low-mass stars, such as 2M1101AB. As they cool down the brown dwarfs should retain a range of luminosities depending on the mass. Without the age and luminosity, a mass estimate is difficult; for example, an L-type brown dwarf could be an old brown dwarf with a high mass (possibly a low-mass star) or a young brown dwarf with a very low mass. For Y dwarfs this is less of a problem, as they remain low-mass objects near the sub-brown dwarf limit, even for relatively high age estimates. For L and T dwarfs it is still useful to have an accurate age estimate. The luminosity is here the less concerning property, as this can be estimated from the spectral energy distribution. The age estimate can be done in two ways. Either the brown dwarf is young and still has spectral features that are associated with youth, or the brown dwarf co-moves with a star or stellar group (star cluster or association), where age estimates are easier to obtain. A very young brown dwarf that was further studied with this method is 2M1207 and the companion 2M1207b. Based on the location, proper motion and spectral signature, this object was determined to belong to the ~8-million-year-old TW Hydrae association, and the mass of the secondary was determined to be 8 ± 2 , below the deuterium burning limit. An example of a very old age obtained by the co-movement method is the brown dwarf + white dwarf binary COCONUTS-1, with the white dwarf estimated to be billion years old. In this case the mass was not estimated with the derived age, but the co-movement provided an accurate distance estimate, using Gaia parallax. Using this measurement the authors estimated the radius, which was then used to estimate the mass for the brown dwarf as .
Observations
Classification of brown dwarfs
Spectral class M
These are brown dwarfs with a spectral class of M5.5 or later; they are also called late-M dwarfs. Some scientists regard them as red dwarfs. All brown dwarfs with spectral type M are young objects, such as Teide 1, which is the first M-type brown dwarf discovered, and LP 944-20, the closest M-type brown dwarf.
Spectral class L
The defining characteristic of spectral class M, the coolest type in the long-standing classical stellar sequence, is an optical spectrum dominated by absorption bands of titanium(II) oxide (TiO) and vanadium(II) oxide (VO) molecules. However, GD 165B, the cool companion to the white dwarf GD 165, had none of the hallmark TiO features of M dwarfs. The subsequent identification of many objects like GD 165B ultimately led to the definition of a new spectral class, the L dwarfs, defined in the red optical region of the spectrum not by metal-oxide absorption bands (TiO, VO), but by metal hydride emission bands (FeH, CrH, MgH, CaH) and prominent atomic lines of alkali metals (Na, K, Rb, Cs). , over 900 L dwarfs had been identified, most by wide-field surveys: the Two Micron All Sky Survey (2MASS), the Deep Near Infrared Survey of the Southern Sky (DENIS), and the Sloan Digital Sky Survey (SDSS). This spectral class also contains the coolest main-sequence stars (> 80 MJ), which have spectral classes L2 to L6.
Spectral class T
As GD 165B is the prototype of the L dwarfs, Gliese 229B is the prototype of a second new spectral class, the T dwarfs. T dwarfs are pinkish-magenta. Whereas near-infrared (NIR) spectra of L dwarfs show strong absorption bands of H2O and carbon monoxide (CO), the NIR spectrum of Gliese 229B is dominated by absorption bands from methane (CH4), a feature which in the Solar System is found only in the giant planets and Titan. CH4, H2O, and molecular hydrogen (H2) collision-induced absorption (CIA) give Gliese 229B blue near-infrared colors. Its steeply sloped red optical spectrum also lacks the FeH and CrH bands that characterize L dwarfs and instead is influenced by exceptionally broad absorption features from the alkali metals Na and K. These differences led J. Davy Kirkpatrick to propose the T spectral class for objects exhibiting H- and K-band CH4 absorption. , 355 T dwarfs were known. NIR classification schemes for T dwarfs have recently been developed by Adam Burgasser and Tom Geballe. Theory suggests that L dwarfs are a mixture of very-low-mass stars and sub-stellar objects (brown dwarfs), whereas the T dwarf class is composed entirely of brown dwarfs. Because of the absorption of sodium and potassium in the green part of the spectrum of T dwarfs, the actual appearance of T dwarfs to human visual perception is estimated to be not brown, but magenta. Early observations limited how distant T-dwarfs could be observed. T-class brown dwarfs, such as WISE 0316+4307, have been detected more than 100 light-years from the Sun. Observations with JWST have detected T-dwarfs such as UNCOVER-BD-1 up to 4500 parsec distant from the sun.
Spectral class Y
In 2009, the coolest-known brown dwarfs had estimated effective temperatures between , and have been assigned the spectral class T9. Three examples are the brown dwarfs CFBDS J005910.90–011401.3, ULAS J133553.45+113005.2 and ULAS J003402.77−005206.7. The spectra of these objects have absorption peaks around 1.55 micrometres. Delorme et al. have suggested that this feature is due to absorption from ammonia and that this should be taken as indicating the T–Y transition, making these objects of type Y0. However, the feature is difficult to distinguish from absorption by water and methane, and other authors have stated that the assignment of class Y0 is premature.
The first JWST spectral energy distribution of a Y-dwarf was able to observe several bands of molecules in the atmosphere of the Y0-dwarf WISE 0359−5401. The observations covered spectroscopy from 1 to 12 μm and photometry at 15, 18 and 21 μm. The molecules water (H2O), methane (CH4), carbon monoxide (CO), carbon dioxide (CO2) and ammonia (NH3) were detected in WISE 0359−5401. Many of these features have been observed before in this Y-dwarf and warmer T-dwarfs by other observatories, but JWST was able to observe them in a single spectrum. Methane is the main reservoir of carbon in the atmosphere of WISE 0359−5401, but there is still enough carbon left to form detectable carbon monoxide (at 4.5–5.0 μm) and carbon dioxide (at 4.2–4.35 μm) in the Y-dwarf. Ammonia was difficult to detect before JWST, as it blends in with the absorption feature of water in the near-infrared, as well at 5.5–7.1 μm. At longer wavelengths of 8.5–12 μm the spectrum of WISE 0359−5401 is dominated by the absorption of ammonia. At 3 μm there is an additional newly detected ammonia feature.
Role of vertical mixing
In the hydrogen-dominated atmosphere of brown dwarfs a chemical equilibrium between carbon monoxide and methane exists. Carbon monoxide reacts with hydrogen molecules and forms methane and hydroxyl in this reaction. The hydroxyl radical might later react with hydrogen and form water molecules. In the other direction of the reaction, methane reacts with hydroxyl and forms carbon monoxide and hydrogen. The chemical reaction is tilted towards carbon monoxide at higher temperatures (L-dwarfs) and lower pressure. At lower temperatures (T-dwarfs) and higher pressure the reaction is tilted towards methane, and methane predominates at the T/Y-boundary. However, vertical mixing of the atmosphere can cause methane to sink into lower layers of the atmosphere and carbon monoxide to rise from these lower and hotter layers. The carbon monoxide is slow to react back into methane because of an energy barrier that prevents the breakdown of the C-O bonds. This forces the observable atmosphere of a brown dwarf to be in a chemical disequilibrium. The L/T transition is mainly defined with the transition from a carbon-monoxide-dominated atmosphere in L-dwarfs to a methane-dominated atmosphere in T-dwarfs. The amount of vertical mixing can therefore push the L/T-transition to lower or higher temperatures. This becomes important for objects with modest surface gravity and extended atmospheres, such as giant exoplanets. This pushes the L/T transition to lower temperatures for giant exoplanets. For brown dwarfs this transition occurs at around 1200 K. The exoplanet HR 8799c, on the other hand, does not show any methane, while having a temperature of 1100K.
The transition between T- and Y-dwarfs is often defined as 500 K because of the lack of spectral observations of these cold and faint objects. Future observations with JWST and the ELTs might improve the sample of Y-dwarfs with observed spectra. Y-dwarfs are dominated by deep spectral features of methane, water vapor and possibly absorption features of ammonia and water ice. Vertical mixing, clouds, metallicity, photochemistry, lightning, impact shocks and metallic catalysts might influence the temperature at which the L/T and T/Y transition occurs.
Secondary features
Young brown dwarfs have low surface gravities because they have larger radii and lower masses than the field stars of similar spectral type. These sources are noted by a letter beta (β) for intermediate surface gravity or gamma (γ) for low surface gravity. Indicators of low surface gravity include weak CaH, K I and Na I lines, as well as a strong VO line. Alpha (α) denotes normal surface gravity and is usually dropped. Sometimes an extremely low surface gravity is denoted by a delta (δ). The suffix "pec" stands for "peculiar"; this suffix is still used for other features that are unusual, and summarizes different properties, indicating low surface gravity, subdwarfs and unresolved binaries. The prefix sd stands for subdwarf and only includes cool subdwarfs. This prefix indicates a low metallicity and kinematic properties that are more similar to halo stars than to disk stars. Subdwarfs appear bluer than disk objects. The red suffix describes objects with red color, but an older age. This is not interpreted as low surface gravity, but as a high dust content. The blue suffix describes objects with blue near-infrared colors that cannot be explained with low metallicity. Some are explained as L+T binaries, others are not binaries, such as 2MASS J11263991−5003550 and are explained with thin and/or large-grained clouds.
Spectral and atmospheric properties of brown dwarfs
The majority of flux emitted by L and T dwarfs is in the 1- to 2.5-micrometre near-infrared range. Low and decreasing temperatures through the late-M, -L, and -T dwarf sequence result in a rich near-infrared spectrum containing a wide variety of features, from relatively narrow lines of neutral atomic species to broad molecular bands, all of which have different dependencies on temperature, gravity, and metallicity. Furthermore, these low temperature conditions favor condensation out of the gas state and the formation of grains.
Typical atmospheres of known brown dwarfs range in temperature from 2200 down to . Compared to stars, which warm themselves with steady internal fusion, brown dwarfs cool quickly over time; more massive dwarfs cool more slowly than less massive ones. There is some evidence that the cooling of brown dwarfs slows down at the transition between spectral classes L and T (about 1000 K).
Observations of known brown dwarf candidates have revealed a pattern of brightening and dimming of infrared emissions that suggests relatively cool, opaque cloud patterns obscuring a hot interior that is stirred by extreme winds. The weather on such bodies is thought to be extremely strong, comparable to but far exceeding Jupiter's famous storms.
On January 8, 2013, astronomers using NASA's Hubble and Spitzer space telescopes probed the stormy atmosphere of a brown dwarf named 2MASS J22282889–4310262, creating the most detailed "weather map" of a brown dwarf thus far. It shows wind-driven, planet-sized clouds. The new research is a stepping stone toward a better understanding not only brown dwarfs, but also of the atmospheres of planets beyond the Solar System.
In April 2020 scientists reported measuring wind speeds of (up to 1,450 miles per hour) on the nearby brown dwarf 2MASS J10475385+2124234. To calculate the measurements, scientists compared the rotational movement of atmospheric features, as ascertained by brightness changes, against the electromagnetic rotation generated by the brown dwarf's interior. The results confirmed previous predictions that brown dwarfs would have high winds. Scientists are hopeful that this comparison method can be used to explore the atmospheric dynamics of other brown dwarfs and extrasolar planets.
Observational techniques
Coronagraphs have recently been used to detect faint objects orbiting bright visible stars, including Gliese 229B.
Sensitive telescopes equipped with charge-coupled devices (CCDs) have been used to search distant star clusters for faint objects, including Teide 1.
Wide-field searches have identified individual faint objects, such as Kelu-1 (30 light-years away).
Brown dwarfs are often discovered in surveys to discover exoplanets. Methods of detecting exoplanets work for brown dwarfs as well, although brown dwarfs are much easier to detect.
Brown dwarfs can be powerful emitters of radio emission due to their strong magnetic fields. Observing programs at the Arecibo Observatory and the Very Large Array have detected over a dozen such objects, which are also called ultracool dwarfs because they share common magnetic properties with other objects in this class. The detection of radio emission from brown dwarfs permits their magnetic field strengths to be measured directly.
Milestones
1995: First brown dwarf verified. Teide 1, an M8 object in the Pleiades cluster, is picked out with a CCD in the Spanish Observatory of Roque de los Muchachos of the Instituto de Astrofísica de Canarias.
First methane brown dwarf verified. Gliese 229B is discovered orbiting red dwarf Gliese 229A (20 ly away) using an adaptive optics coronagraph to sharpen images from the reflecting telescope at Palomar Observatory on Southern California's Mount Palomar; follow-up infrared spectroscopy made with their Hale Telescope shows an abundance of methane.
1998: First X-ray-emitting brown dwarf found. Cha Helpha 1, an M8 object in the Chamaeleon I dark cloud, is determined to be an X-ray source, similar to convective late-type stars.
15 December 1999: First X-ray flare detected from a brown dwarf. A team at the University of California monitoring LP 944-20 (, 16 ly away) via the Chandra X-ray Observatory, catches a 2-hour flare.
27 July 2000: First radio emission (in flare and quiescence) detected from a brown dwarf. A team of students at the Very Large Array detected emission from LP 944–20.
30 April 2004: First detection of a candidate exoplanet around a brown dwarf: 2M1207b discovered with the VLT and the first directly imaged exoplanet.
20 March 2013: Discovery of the closest brown dwarf system: Luhman 16.
25 April 2014: Coldest-known brown dwarf discovered. WISE 0855−0714 is 7.2 light-years away (seventh-closest system to the Sun) and has a temperature between −48 and −13 °C.
Brown dwarfs X-ray sources
X-ray flares detected from brown dwarfs since 1999 suggest changing magnetic fields within them, similar to those in very-low-mass stars. Although they do not fuse hydrogen into helium in their cores like stars, energy from the fusion of deuterium and gravitational contraction keep their interiors warm and generate strong magnetic fields. The interior of a brown dwarf is in a rapidly boiling, or convective state. When combined with the rapid rotation that most brown dwarfs exhibit, convection sets up conditions for the development of a strong, tangled magnetic field near the surface. The magnetic fields that generated the flare observed by Chandra from LP 944-20 has its origin in the turbulent magnetized plasma beneath the brown dwarf's "surface".
Using NASA's Chandra X-ray Observatory, scientists have detected X-rays from a low-mass brown dwarf in a multiple star system. This is the first time that a brown dwarf this close to its parent star(s) (Sun-like stars TWA 5A) has been resolved in X-rays. "Our Chandra data show that the X-rays originate from the brown dwarf's coronal plasma which is some 3 million degrees Celsius", said Yohko Tsuboi of Chuo University in Tokyo. "This brown dwarf is as bright as the Sun today in X-ray light, while it is fifty times less massive than the Sun", said Tsuboi. "This observation, thus, raises the possibility that even massive planets might emit X-rays by themselves during their youth!"
Brown dwarfs as radio sources
The first brown dwarf that was discovered to emit radio signals was LP 944-20, which was observed since it is also a source of X-ray emission, and both types of emission are signatures of coronae. Approximately 5–10% of brown dwarfs appear to have strong magnetic fields and emit radio waves, and there may be as many as 40 magnetic brown dwarfs within 25 pc of the Sun based on Monte Carlo modeling and their average spatial density. The power of the radio emissions of brown dwarfs is roughly constant despite variations in their temperatures. Brown dwarfs may maintain magnetic fields of up to 6 kG in strength. Astronomers have estimated brown dwarf magnetospheres to span an altitude of approximately 107 m given properties of their radio emissions. It is unknown whether the radio emissions from brown dwarfs more closely resemble those from planets or stars. Some brown dwarfs emit regular radio pulses, which are sometimes interpreted as radio emission beamed from the poles but may also be beamed from active regions. The regular, periodic reversal of radio wave orientation may indicate that brown dwarf magnetic fields periodically reverse polarity. These reversals may be the result of a brown dwarf magnetic activity cycle, similar to the solar cycle.
The first brown dwarf of spectral class M found to emit radio waves was LP 944-20, detected in 2001. The first brown dwarf of spectral class L found to emit radio waves was 2MASS J0036159+182110, detected in 2008. The first brown dwarf of spectral class T found to emit radio waves was 2MASS J10475385+2124234. This last discovery was significant since it revealed that brown dwarfs with temperatures similar to exoplanets could host strong >1.7 kG magnetic fields. Although a sensitive search for radio emission from Y dwarfs was conducted at the Arecibo Observatory in 2010, no emission was detected.
Recent developments
Estimates of brown dwarf populations in the solar neighbourhood suggest that there may be as many as six stars for every brown dwarf. A more recent estimate from 2017 using the young massive star cluster RCW 38 concluded that the Milky Way galaxy contains between 25 and 100 billion brown dwarfs. (Compare these numbers to the estimates of the number of stars in the Milky Way; 100 to 400 billion.)
In a study published in Aug 2017 NASA's Spitzer Space Telescope monitored infrared brightness variations in brown dwarfs caused by cloud cover of variable thickness. The observations revealed large-scale waves propagating in the atmospheres of brown dwarfs (similarly to the atmosphere of Neptune and other Solar System giant planets). These atmospheric waves modulate the thickness of the clouds and propagate with different velocities (probably due to differential rotation).
In August 2020, astronomers discovered 95 brown dwarfs near the Sun through the project Backyard Worlds: Planet 9.
In 2024 the James Webb Space Telescope provided the most detailed weather report yet on two brown dwarfs, revealing "stormy" conditions. These brown dwarfs, part of a binary star system named Luhman 16 discovered in 2013, are only 6.5 light-years away from Earth and are the closest brown dwarfs to our sun. Researchers discovered that they have turbulent clouds, likely made of silicate grains, with temperatures ranging from to . This indicates that hot sand is being blown by winds on the brown dwarfs. Additionally, absorption signatures of carbon monoxide, methane, and water vapor were detected.
Binary brown dwarfs
Brown dwarf–brown dwarf binaries
Brown dwarfs binaries of type M, L, and T are less common with a lower mass of the primary. L-dwarfs have a binary fraction of about % and the binary fraction for late T, early Y-dwarfs (T5-Y0) is about .
Brown dwarf binaries have a higher companion-to-host ratio for lower mass binaries. Binaries with a M-type star as a primary have for example a broad distribution of q with a preference of q ≥ 0.4. Brown dwarfs on the other hand show a strong preference for q ≥ 0.7. The separation is decreasing with mass: M-type stars have a separation peaking at 3–30 astronomical units (au), M-L-type brown dwarfs have a projected separation peaking at 5–8 au and T5–Y0 objects have a projected separation that follows a lognormal distribution with a peak separation of about 2.9 au.
An example is the closest brown dwarf binary Luhman 16 AB with a primary L7.5 dwarf and a separation of 3.5 au and q = 0.85. The separation is on the lower end of the expected separation for M-L-type brown dwarfs, but the mass ratio is typical.
It is not known if the same trend continues with Y-dwarfs, because their sample size is so small. The Y+Y dwarf binaries should have a high mass ratio q and a low separation, reaching scales of less than one au. In 2023, the Y+Y dwarf WISE J0336-0143 was confirmed as a binary with JWST, with a mass ratio of q=0.62±0.05 and a separation of 0.97 astronomical units. The researchers point out that the sample size of low-mass binary brown dwarfs is too small to determine if WISE J0336-0143 is a typical representative of low-mass binaries or a peculiar system.
Observations of the orbit of binary systems containing brown dwarfs can be used to measure the mass of the brown dwarf. In the case of 2MASSW J0746425+2000321, the secondary weighs 6% of the solar mass. This measurement is called a dynamical mass. The brown dwarf system closest to the Solar System is the binary Luhman 16. It was attempted to search for planets around this system with a similar method, but none were found.
Unusual brown dwarf binaries
The wide binary system 2M1101AB was the first binary with a separation greater than . The discovery of the system gave definitive insights to the formation of brown dwarfs. It was previously thought that wide binary brown dwarfs are not formed or at least are disrupted at ages of 1–10 Myr. The existence of this system is also inconsistent with the ejection hypothesis. The ejection hypothesis was a proposed hypothesis in which brown dwarfs form in a multiple system, but are ejected before they gain enough mass to burn hydrogen.
More recently the wide binary W2150AB was discovered. It has a similar mass ratio and binding energy as 2M1101AB, but a greater age and is located in a different region of the galaxy. While 2M1101AB is in a closely crowded region, the binary W2150AB is in a sparsely-separated field. It must have survived any dynamical interactions in its natal star cluster. The binary belongs also to a few L+T binaries that can be easily resolved by ground-based observatories. The other two are SDSS J1416+13AB and Luhman 16.
There are other interesting binary systems such as the eclipsing binary brown dwarf system 2MASS J05352184–0546085. Photometric studies of this system have revealed that the less massive brown dwarf in the system is hotter than its higher-mass companion.
Brown dwarfs around stars
Brown dwarfs and massive planets in a close orbit (less than 5 au) around stars are rare and this is sometimes described as the brown dwarf desert. Less than 1% of stars with the mass of the sun have a brown dwarf within 3–5 au.
An example for a star–brown dwarf binary is the first discovered T-dwarf Gliese 229 B, which orbits around the main-sequence star Gliese 229 A, a red dwarf. Brown dwarfs orbiting subgiants are also known, such as TOI-1994b which orbits its star every 4.03 days.
There is also disagreement if some low-mass brown dwarfs should be considered planets. The NASA Exoplanet archive includes brown dwarfs with a minimum mass less or equal to 30 Jupiter masses as planets as long as there are other criteria fulfilled (e.g. orbiting a star). The Working Group on Extrasolar Planets (WGESP) of the IAU on the other hand only considers planets with a mass below 13 Jupiter masses.
White dwarf–brown dwarf binaries
Brown dwarfs around white dwarfs are quite rare. GD 165 B, the prototype of the L dwarfs, is one such system. Such systems can be useful in determining the age of the system and the mass of the brown dwarf. Other white dwarf-brown dwarf binaries are COCONUTS-1 AB (7 billion years old), and LSPM J0055+5948 AB (10 billion years old), SDSS J22255+0016 AB (2 billion years old) WD 0806−661 AB (1.5–2.7 billion years old).
Systems with close, tidally locked brown dwarfs orbiting around white dwarfs belong to the post common envelope binaries or PCEBs. Only eight confirmed PCEBs containing a white dwarf with a brown dwarf companion are known, including WD 0137-349 AB. In the past history of these close white dwarf–brown dwarf binaries, the brown dwarf is engulfed by the star in the red giant phase. Brown dwarfs with a mass lower than 20 Jupiter masses would evaporate during the engulfment. The dearth of brown dwarfs orbiting close to white dwarfs can be compared with similar observations of brown dwarfs around main-sequence stars, described as the brown-dwarf desert. The PCEB might evolve into a cataclysmic variable star (CV*) with the brown dwarf as the donor. Simulations have shown that highly evolved CV* are mostly associated with substellar donors (up to 80%). A type of CV*, called WZ Sge-type dwarf nova often show donors with a mass near the borderline of low-mass stars and brown dwarfs. The binary BW Sculptoris is such a dwarf nova with a brown dwarf donor. This brown dwarf likely formed when a donor star lost enough mass to become a brown dwarf. The mass loss comes with a loss of the orbital period until it reaches a minimum of 70–80 minutes at which the period increases again. This gives this evolutionary stage the name period bouncer. There could also exist brown dwarfs that merged with white dwarfs. The nova CK Vulpeculae might be a result of such a white dwarf–brown dwarf merger.
Formation and evolution
The earliest stage of brown dwarf formation is called proto- or pre-brown dwarf. Proto-brown dwarfs are low-mass equivalents of protostars (class 0/I objects). Additionally Very Low Luminosity Objects (VeLLOs) that have Lint ≤0.1-0.2 are often proto-brown dwarfs. They are found in nearby star-forming clouds. Around 67 promising proto-brown dwarfs and 26 pre-brown dwarfs are known as of 2024. As of 2017 there is only one known proto-brown dwarf that is connected with a large Herbig–Haro object. This is the brown dwarf Mayrit 1701117, which is surrounded by a pseudo-disk and a Keplerian disk. Mayrit 1701117 launches the 0.7-light-year-long jet HH 1165, mostly seen in ionized sulfur.
Brown dwarfs form similarly to stars and are surrounded by protoplanetary disks, such as Cha 110913−773444. Disks around brown dwarfs have been found to have many of the same features as disks around stars; therefore, it is expected that there will be accretion-formed planets around brown dwarfs. Given the small mass of brown dwarf disks, most planets will be terrestrial planets rather than gas giants. If a giant planet orbits a brown dwarf across our line of sight, then, because they have approximately the same diameter, this would give a large signal for detection by transit. The accretion zone for planets around a brown dwarf is very close to the brown dwarf itself, so tidal forces would have a strong effect.
In 2020, the closest brown dwarf with an associated primordial disk (class II disk)—WISEA J120037.79-784508.3 (W1200-7845)—was discovered by the Disk Detective project when classification volunteers noted its infrared excess. It was vetted and analyzed by the science team who found that W1200-7845 had a 99.8% probability of being a member of the ε Chamaeleontis (ε Cha) young moving group association. Its parallax (using Gaia DR2 data) puts it at a distance of 102 parsecs (or 333 lightyears) from Earth—which is within the local Solar neighborhood.
A paper from 2021 studied circumstellar discs around brown dwarfs in stellar associations that are a few million years old and 140 to 200 parsecs away. The researchers found that these disks are not massive enough to form planets in the future. There is evidence in these disks that might indicate that planet formation begins at earlier stages and that planets are already present in these disks. The evidence for disk evolution includes a decreasing disk mass over time, dust grain growth and dust settling. Two brown dwarf disks were also found in absorption and at least 4 disks are photoevaporating from external UV-ratiation in the Orion Nebula. Such objects are also called proplyds. Proplyd 181−247, which is a brown dwarf or low-mass star, is surrounded by a disk with a radius of 30 astronomical units and the disk has a mass of 6.2±1.0 . Disks around brown dwarfs usually have a radius smaller than 40 astronomical units, but three disks in the more distant Taurus molecular cloud have a radius larger than 70 au and were resolved with ALMA. These larger disks are able to form rocky planets with a mass >1 . There are also brown dwarfs with disks in associations older than a few million years, which might be evidence that disks around brown dwarfs need more time to dissipate. Especially old disks (>20 Myrs) are sometimes called Peter Pan disks. Currently 2MASS J02265658-5327032 is the only known brown dwarf that has a Peter Pan disk.
The brown dwarf Cha 110913−773444, located 500 light-years away in the constellation Chamaeleon, may be in the process of forming a miniature planetary system. Astronomers from Pennsylvania State University have detected what they believe to be a disk of gas and dust similar to the one hypothesized to have formed the Solar System. Cha 110913−773444 is the smallest brown dwarf found to date (), and if it formed a planetary system, it would be the smallest-known object to have one.
Planets around brown dwarfs
According to the IAU working definition (from August 2018) an exoplanet can orbit a brown dwarf. It requires a mass below 13 and a mass ratio of M/Mcentral<2/(25+√621). This means that an object with a mass up to 3.2 around a brown dwarf with a mass of 80 is considered a planet. It also means that an object with a mass up to 0.52 around a brown dwarf with a mass of 13 is considered a planet.
The super-Jupiter planetary-mass objects 2M1207b, 2MASS J044144 and Oph 98 B that are orbiting brown dwarfs at large orbital distances may have formed by cloud collapse rather than accretion and so may be sub-brown dwarfs rather than planets, which is inferred from relatively large masses and large orbits. The first discovery of a low-mass companion orbiting a brown dwarf (ChaHα8) at a small orbital distance using the radial velocity technique paved the way for the detection of planets around brown dwarfs on orbits of a few AU or smaller. However, with a mass ratio between the companion and primary in ChaHα8 of about 0.3, this system rather resembles a binary star. Then, in 2008, the first planetary-mass companion in a relatively small orbit (MOA-2007-BLG-192Lb) was discovered orbiting a brown dwarf.
Planets around brown dwarfs are likely to be carbon planets depleted of water.
A 2017 study, based upon observations with Spitzer estimates that 175 brown dwarfs need to be monitored in order to guarantee (95%) at least one detection of a below earth-sized planet via the transiting method. JWST could potentially detect smaller planets. The orbits of planets and moons in the solar system often align with the orientation of the host star/planet they orbit. Assuming the orbit of a planet is aligned with the rotational axis of a brown dwarf or planetary-mass object, the geometric transit probability of an object similar to Io can be calculated with the formula cos(79.5°)/cos(inclination). The inclination was estimated for several brown dwarfs and planetary-mass objects. SIMP 0136 for example has an estimated inclination of 80°±12. Assuming the lower bound of i≥68° for SIMP 0136, this results in a transit probability of ≥48.6% for close-in planets. It is however not known how common close-in planets are around brown dwarfs and they might be more common for lower-mass objects, as disk sizes seem to decrease with mass.
Habitability
Habitability for hypothetical planets orbiting brown dwarfs has been studied. Computer models suggesting conditions for these bodies to have habitable planets are very stringent, the habitable zone being narrow, close (T dwarf 0.005 au) and decreasing with time, due to the cooling of the brown dwarf (they fuse for at most 10 million years). The orbits there would have to be of extremely low eccentricity (on the order of 10 to the minus 6) to avoid strong tidal forces that would trigger a runaway greenhouse effect on the planets, rendering them uninhabitable. There would also be no moons.
Superlative brown dwarfs
In 1984, it was postulated by some astronomers that the Sun may be orbited by an undetected brown dwarf (sometimes referred to as Nemesis) that could interact with the Oort cloud just as passing stars can. However, this hypothesis has fallen out of favor.
Table of firsts
Table of extremes
| Physical sciences | Stellar astronomy | null |
44410 | https://en.wikipedia.org/wiki/Conjunctivitis | Conjunctivitis | Conjunctivitis, also known as pink eye or Madras eye, is inflammation of the conjunctiva and the inner surface of the eyelid. It makes the eye appear pink or reddish. Pain, burning, scratchiness, or itchiness may occur. The affected eye may have increased tears or be "stuck shut" in the morning. Swelling of the sclera may also occur. Itching is more common in cases due to allergies. Conjunctivitis can affect one or both eyes.
The most common infectious causes in adults are viral, whereas in children bacterial causes predominate. The viral infection may occur along with other symptoms of a common cold. Both viral and bacterial cases are easily spread between people. Allergies to pollen or animal hair are also a common cause. Diagnosis is often based on signs and symptoms. Occasionally, a sample of the discharge is sent for culture.
Prevention is partly by handwashing. Treatment depends on the underlying cause. In the majority of viral cases, there is no specific treatment. Most cases due to a bacterial infection also resolve without treatment; however, antibiotics can shorten the illness. People who wear contact lenses and those whose infection is caused by gonorrhea or chlamydia should be treated. Allergic cases can be treated with antihistamines or mast cell inhibitor drops.
About 3 to 6 million people get acute conjunctivitis each year in the United States. Typically, people get better in one or two weeks. If visual loss, significant pain, sensitivity to light or signs of herpes occur, or if symptoms do not improve after a week, further diagnosis and treatment may be required. Conjunctivitis in a newborn, known as neonatal conjunctivitis, may also require specific treatment.
Signs and symptoms
Red eye, swelling of the conjunctiva, and watering of the eyes are symptoms common to all forms of conjunctivitis. However, the pupils should be normally reactive, and the visual acuity normal.
Conjunctivitis is identified by inflammation of the conjunctiva, manifested by irritation and redness. Examination using a slit lamp (biomicroscope) may improve diagnostic accuracy. Examination of the palpebral conjunctiva, that overlying the inner aspects of the eyelids, is usually more diagnostic than examination of the bulbal conjunctiva, that overlying the sclera.
Viral
Approximately 80% of cases of conjunctivitis in adults and less than 20% in children are due to viruses, with 65% to 90% of these cases being attributed to adenoviruses.
Viral conjunctivitis is often associated with an infection of the upper respiratory tract, a common cold, or a sore throat. Other associated signs may include pre-auricular lymph node swelling and contact with another person with a red eye. Eye pain may be present if the cornea is also involved. Its symptoms include excessive watering and itching. The discharge in viral conjunctivitis is usually (but not always) watery in nature. The infection usually begins in one eye but may spread easily to the other eye.
Viral conjunctivitis manifests as a fine, diffuse pinkness of the conjunctiva which may be mistaken for iritis, but corroborative signs on microscopy, particularly numerous lymphoid follicles on the tarsal conjunctiva, and sometimes a punctate keratitis are seen.
Allergic
Allergic conjunctivitis is inflammation of the conjunctiva due to allergy. The specific allergens may differ among patients. Symptoms result from the release of histamine and other active substances by mast cells, and consist of redness (mainly due to vasodilation of the peripheral small blood vessels), swelling of the conjunctiva, itching, and increased production of tears.
Bacterial
Bacteria are responsible for approximately 70% of conjunctivitis in children and less than 20% of cases in adults. Common bacteria responsible for bacterial conjunctivitis are Staphylococcus including Staph aureus, Streptococcus such as strep pneumoniae, Haemophilus species and Moraxella catarrhalis. Less commonly, Chlamydia spp. and Niesseria species (Neisseria gonorrhoeae and Neisseria meningitidis) may be the cause. Infection with Escherichia coli may also cause conjunctivitis, particularly in the neonatal subtype ophthalmia neonatorum. Bacterial conjunctivitis usually causes a rapid onset of conjunctival redness, swelling of the eyelid, and a sticky discharge. Typically, symptoms develop first in one eye, but may spread to the other eye within 2–5 days. Conjunctivitis due to common pus-producing bacteria causes marked grittiness or irritation and a stringy, opaque, greyish or yellowish discharge that may cause the lids to stick together, especially after sleep. Severe crusting of the infected eye and the surrounding skin may also occur. The gritty or scratchy feeling is sometimes localized enough that patients may insist that they have a foreign body in the eye.
Bacteria such as Chlamydia trachomatis or Moraxella spp. can cause a nonexudative but persistent conjunctivitis without much redness. Bacterial conjunctivitis may cause the production of membranes or pseudomembranes that cover the conjunctiva. Pseudomembranes consist of a combination of inflammatory cells and exudates and adhere loosely to the conjunctiva, while true membranes are more tightly adherent and cannot be easily peeled away. Cases of bacterial conjunctivitis that involve the production of membranes or pseudomembranes are associated with Neisseria gonorrhoeae, β-hemolytic streptococci, and Corynebacterium diphtheriae. C. diphtheriae causes membrane formation in conjunctiva of unimmunized children.
Chemical
Chemical eye injury may result when an acidic or alkaline substance gets in the eye. Alkali burns are typically worse than acidic burns. Mild burns produce conjunctivitis, while more severe burns may cause the cornea to turn white. Litmus paper may be used to test for chemical causes. When a chemical cause has been confirmed, the eye or eyes should be flushed until the pH is in the range 6–8. Anaesthetic eye drops can be used to decrease the pain.
Irritant or toxic conjunctivitis is primarily marked by redness. If due to a chemical splash, it is often present in only the lower conjunctival sac. With some chemicals, above all with caustic alkalis such as sodium hydroxide, necrosis of the conjunctiva marked by a deceptively white eye due to vascular closure may occur, followed by sloughing off of the dead epithelium. A slit lamp examination is likely to show evidence of anterior uveitis.
Biomarkers
Omics technologies have been used to identify biomarkers that inform on the emergence and progression of conjunctivitis. For example, in chronic inflammatory cicatrizing conjunctivitis, active oxylipins, lysophospholipids, fatty acids, and endocannabinoids alterations, from which potential biomarkers linked to inflammatory processes were identified.
Other
Inclusion conjunctivitis of the newborn is a conjunctivitis that may be caused by the bacterium Chlamydia trachomatis, and may lead to acute, purulent conjunctivitis. However, it is usually self-healing.
Causes
Infective conjunctivitis is most commonly caused by a virus. Bacterial infections, allergies, other irritants, and dryness are also common causes. Both bacterial and viral infections are contagious, passing from person to person or spread through contaminated objects or water. Contact with contaminated fingers is a common cause of conjunctivitis. Bacteria may also reach the conjunctiva from the edges of the eyelids and the surrounding skin, from the nasopharynx, from infected eye drops or contact lenses, from the genitals or the bloodstream. Infection by human adenovirus accounts for 65% to 90% of cases of viral conjunctivitis.
Viral
Adenoviruses are the most common cause of viral conjunctivitis (adenoviral keratoconjunctivitis). Herpetic keratoconjunctivitis, caused by herpes simplex viruses, can be serious and requires treatment with aciclovir. Acute hemorrhagic conjunctivitis is a highly contagious disease caused by one of two enteroviruses, enterovirus 70 and coxsackievirus A24. These were first identified in an outbreak in Ghana in 1969, and have spread worldwide since then, causing several epidemics.
Bacterial
The most common causes of acute bacterial conjunctivitis are Staphylococcus aureus, Streptococcus pneumoniae, and Haemophilus influenzae. Though very rare, hyperacute cases are usually caused by Neisseria gonorrhoeae or Neisseria meningitidis. Chronic cases of bacterial conjunctivitis are those lasting longer than 3 weeks, and are typically caused by S. aureus, Moraxella lacunata, or Gram-negative enteric flora.
Allergic
Conjunctivitis may also be caused by allergens such as pollen, perfumes, cosmetics, smoke, dust mites, Balsam of Peru, or eye drops. The most frequent cause of conjunctivitis is allergic conjunctivitis and it affects 15% to 40% of the population. Allergic conjunctivitis accounts for 15% of eye related primary care consultations; most including seasonal exposures in the spring and summer or perpetual conditions.
Other
Computer vision syndrome
Dry eye syndrome
Reactive arthritis: Conjunctivitis is part of the triad of reactive arthritis, which is thought to be caused by autoimmune cross-reactivity following certain bacterial infections. Reactive arthritis is highly associated with HLA-B27. Conjunctivitis is associated with the autoimmune disease relapsing polychondritis.
Diagnosis
Cultures are not often taken or needed as most cases resolve either with time or typical antibiotics. If bacterial conjunctivitis is suspected, but no response to topical antibiotics is seen, swabs for bacterial culture should be taken and tested. Viral culture may be appropriate in epidemic case clusters.
A patch test is used to identify the causative allergen in allergic conjunctivitis.
Although conjunctival scrapes for cytology can be useful in detecting chlamydial and fungal infections, allergies, and dysplasia, they are rarely done because of the cost and the general dearth of laboratory staff experienced in handling ocular specimens. Conjunctival incisional biopsy is occasionally done when granulomatous diseases (e.g., sarcoidosis) or dysplasia are suspected.
Classification
Conjunctivitis may be classified either by cause or by extent of the inflamed area.
Causes
Allergy
Bacteria
Viruses
Chemicals
Autoimmune
Neonatal conjunctivitis is often grouped separately from bacterial conjunctivitis because it is caused by different bacteria than the more common cases of bacterial conjunctivitis.
By extent of involvement
Blepharoconjunctivitis is the dual combination of conjunctivitis with blepharitis (inflammation of the eyelids).
Keratoconjunctivitis is the combination of conjunctivitis and keratitis (corneal inflammation).
Blepharokeratoconjunctivitis is the combination of conjunctivitis with blepharitis and keratitis. It is clinically defined by changes of the lid margin, meibomian gland dysfunction, redness of the eye, conjunctival chemosis and inflammation of the cornea.
Differential diagnosis
Some more serious conditions can present with a red eye, such as infectious keratitis, angle-closure glaucoma, or iritis. These conditions require the urgent attention of an ophthalmologist. Signs of such conditions include decreased vision, significantly increased sensitivity to light, inability to keep the eye open, a pupil that does not respond to light, or a severe headache with nausea. Fluctuating blurring is common, due to tearing and mucoid discharge. Mild photophobia is common. However, if any of these symptoms is prominent, considering other diseases such as glaucoma, uveitis, keratitis, and even meningitis or carotico-cavernous fistula is important.
A more comprehensive differential diagnosis for the red or painful eye includes:
Corneal abrasion
Subconjunctival hemorrhage
Pinguecula
Blepharitis
Dacryocystitis
Keratoconjunctivitis sicca (dry eye)
Keratitis
Herpes simplex
Herpes zoster
Episcleritis – an inflammatory condition that produces a similar appearance to conjunctivitis, but without discharge or tearing
Uveitis
Acute angle-closure glaucoma
Endophthalmitis
Orbital cellulitis
Prevention
The most effective prevention is good hygiene, especially avoiding rubbing the eyes with infected hands. Vaccination against some of the causative pathogens such as Haemophilus influenzae, pneumococcus, and Neisseria meningitidis is also effective.
Povidone-iodine eye solution has been found to prevent neonatal conjunctivitis. It is becoming more commonly used globally because of its low cost.
Management
Conjunctivitis resolves in 65% of cases without treatment, within 2–5 days. The prescription of antibiotics is not necessary in most cases.
Viral
Viral conjunctivitis usually resolves on its own and does not require any specific treatment. Antihistamines (e.g., diphenhydramine) or mast cell stabilizers (e.g., cromolyn) may be used to help with the symptoms. Povidone-iodine has been suggested as a treatment, but as of 2008, evidence to support it was poor.
Allergic
For allergic conjunctivitis, cool water poured over the face with the head inclined downward constricts capillaries, and artificial tears sometimes relieve discomfort in mild cases. In more severe cases, nonsteroidal anti-inflammatory medications and antihistamines may be prescribed. Persistent allergic conjunctivitis may also require topical steroid drops.
Bacterial
Bacterial conjunctivitis usually resolves without treatment. Topical antibiotics may be needed only if no improvement is observed after 3 days. No serious effects were noted either with or without treatment. Because antibiotics do speed healing in bacterial conjunctivitis, their use may be considered. Antibiotics are also recommended for those who wear contact lenses, are immunocompromised, have disease which is thought to be due to chlamydia or gonorrhea, have a fair bit of pain, or have copious discharge. Gonorrheal or chlamydial infections require both oral and topical antibiotics.
The choice of antibiotic varies based on the strain or suspected strain of bacteria causing the infection. Fluoroquinolones, sodium sulfacetamide, or trimethoprim/polymyxin may be used, typically for 7–10 days. Cases of meningococcal conjunctivitis can also be treated with systemic penicillin, as long as the strain is sensitive to penicillin.
When investigated as a treatment, povidone-iodine ophthalmic solution has also been observed to have some effectiveness against bacterial and chlamydial conjunctivitis, with a possible role suggested in locations where topical antibiotics are unavailable or costly.
Chemical
Conjunctivitis due to chemicals is treated via irrigation with Ringer's lactate or saline solution. Chemical injuries, particularly alkali burns, are medical emergencies, as they can lead to severe scarring and intraocular damage. People with chemically induced conjunctivitis should not touch their eyes to avoid spreading the chemical.
Epidemiology
Conjunctivitis is the most common eye disease. Rates of disease is related to the underlying cause which varies by the age as well as the time of year. Acute conjunctivitis is most frequently found in infants, school-age children and the elderly. The most common cause of infectious conjunctivitis is viral conjunctivitis.
It is estimated that acute conjunctivitis affects 6 million people annually in the United States.
Some seasonal trends have been observed for the occurrence of different forms of conjunctivitis. In the northern hemisphere, the occurrence of bacterial conjunctivitis peaks from December to April, viral conjunctivitis peaks in the summer months and allergic conjunctivitis is more prevalent throughout the spring and summer.
History
An adenovirus was first isolated by Rowe et al. in 1953. Two years later, Jawetz et al. published on epidemic keratoconjunctivitis. "Madras eye" is a colloquial term that has been used in India for the disease.
Outbreak in Pakistan
In September 2023, a significant outbreak of conjunctivitis occurred in Pakistan. The outbreak began in Karachi and quickly spread to Lahore, Rawalpindi, and Islamabad. By the end of the month, over 86,133 cases had been reported in Punjab alone. The rapid spread of the disease led to the temporary closure of schools in the region. This event marked one of the largest outbreaks of Pink Eye in the country's recent history.
Society and culture
Conjunctivitis imposes economic and social burdens. The cost of treating bacterial conjunctivitis in the United States was estimated to be $377 million to $857 million per year. Approximately 1% of all primary care office visits in the United States are related to conjunctivitis. Approximately 70% of all people with acute conjunctivitis present to primary care and urgent care.
| Biology and health sciences | Infectious diseases by site | Health |
44412 | https://en.wikipedia.org/wiki/Sedimentary%20rock | Sedimentary rock | Sedimentary rocks are types of rock that are formed by the accumulation or deposition of sediments, ie. mineral or organic particles, at Earth's surface, followed by cementation. Sedimentation is the collective name for processes that cause these particles to settle in place. The particles that form a sedimentary rock are called sediment, and may be composed of geological detritus (minerals) or biological detritus (organic matter). The geological detritus originated from weathering and erosion of existing rocks, or from the solidification of molten lava blobs erupted by volcanoes. The geological detritus is transported to the place of deposition by water, wind, ice or mass movement, which are called agents of denudation. Biological detritus was formed by bodies and parts (mainly shells) of dead aquatic organisms, as well as their fecal mass, suspended in water and slowly piling up on the floor of water bodies (marine snow). Sedimentation may also occur as dissolved minerals precipitate from water solution.
The sedimentary rock cover of the continents of the Earth's crust is extensive (73% of the Earth's current land surface), but sedimentary rock is estimated to be only 8% of the volume of the crust. Sedimentary rocks are only a thin veneer over a crust consisting mainly of igneous and metamorphic rocks. Sedimentary rocks are deposited in layers as strata, forming a structure called bedding. Sedimentary rocks are often deposited in large structures called sedimentary basins. Sedimentary rocks have also been found on Mars.
The study of sedimentary rocks and rock strata provides information about the subsurface that is useful for civil engineering, for example in the construction of roads, houses, tunnels, canals or other structures. Sedimentary rocks are also important sources of natural resources including coal, fossil fuels, drinking water and ores.
The study of the sequence of sedimentary rock strata is the main source for an understanding of the Earth's history, including palaeogeography, paleoclimatology and the history of life. The scientific discipline that studies the properties and origin of sedimentary rocks is called sedimentology. Sedimentology is part of both geology and physical geography and overlaps partly with other disciplines in the Earth sciences, such as pedology, geomorphology, geochemistry and structural geology.
Classification based on origin
Sedimentary rocks can be subdivided into four groups based on the processes responsible for their formation: clastic sedimentary rocks, biochemical (biogenic) sedimentary rocks, chemical sedimentary rocks, and a fourth category for "other" sedimentary rocks formed by impacts, volcanism, and other minor processes.
Clastic sedimentary rocks
Clastic sedimentary rocks are composed of rock fragments (clasts) that have been cemented together. The clasts are commonly individual grains of quartz, feldspar, clay minerals, or mica. However, any type of mineral may be present. Clasts may also be lithic fragments composed of more than one mineral.
Clastic sedimentary rocks are subdivided according to the dominant particle size. Most geologists use the Udden-Wentworth grain size scale and divide unconsolidated sediment into three fractions: gravel (>2 mm diameter), sand (1/16 to 2 mm diameter), and mud (<1/16 mm diameter). Mud is further divided into silt (1/16 to 1/256 mm diameter) and clay (<1/256 mm diameter). The classification of clastic sedimentary rocks parallels this scheme; conglomerates and breccias are made mostly of gravel, sandstones are made mostly of sand, and mudrocks are made mostly of mud. This tripartite subdivision is mirrored by the broad categories of rudites, arenites, and lutites, respectively, in older literature.
The subdivision of these three broad categories is based on differences in clast shape (conglomerates and breccias), composition (sandstones), or grain size or texture (mudrocks).
Conglomerates and breccias
Breccias are dominantly composed of angular gravel in a groundmass (matrix), while conglomerates are dominantly composed rounded gravel.
Sandstones
Sandstone classification schemes vary widely, but most geologists have adopted the Dott scheme, which uses the relative abundance of quartz, feldspar, and lithic framework grains and the abundance of a muddy matrix between the larger grains.
Composition of framework grains
The relative abundance of sand-sized framework grains determines the first word in a sandstone name. Naming depends on the dominance of the three most abundant components quartz, feldspar, or the lithic fragments that originated from other rocks. All other minerals are considered accessories and not used in the naming of the rock, regardless of abundance.
Quartz sandstones have >90% quartz grains
Feldspathic sandstones have <90% quartz grains and more feldspar grains than lithic grains
Lithic sandstones have <90% quartz grains and more lithic grains than feldspar grains
Abundance of muddy matrix material between sand grains
When sand-sized particles are deposited, the space between the grains either remains open or is filled with mud (silt and/or clay sized particle).
"Clean" sandstones with open pore space (that may later be filled with matrix material) are called arenites.
Muddy sandstones with abundant (>10%) muddy matrix are called wackes.
Six sandstone names are possible using the descriptors for grain composition (quartz-, feldspathic-, and lithic-) and the amount of matrix (wacke or arenite). For example, a quartz arenite would be composed of mostly (>90%) quartz grains and have little or no clayey matrix between the grains, a lithic wacke would have abundant lithic grains and abundant muddy matrix, etc.
Although the Dott classification scheme is widely used by sedimentologists, common names like greywacke, arkose, and quartz sandstone are still widely used by non-specialists and in popular literature.
Mudrocks
Mudrocks are sedimentary rocks composed of at least 50% silt- and clay-sized particles. These relatively fine-grained particles are commonly transported by turbulent flow in water or air, and deposited as the flow calms and the particles settle out of suspension.
Most authors presently use the term "mudrock" to refer to all rocks composed dominantly of mud. Mudrocks can be divided into siltstones, composed dominantly of silt-sized particles; mudstones with subequal mixture of silt- and clay-sized particles; and claystones, composed mostly of clay-sized particles. Most authors use "shale" as a term for a fissile mudrock (regardless of grain size) although some older literature uses the term "shale" as a synonym for mudrock.
Biochemical sedimentary rocks
Biochemical sedimentary rocks are created when organisms use materials dissolved in air or water to build their tissue. Examples include:
Most types of limestone are formed from the calcareous skeletons of organisms such as corals, mollusks, and foraminifera.
Coal, formed from vegetation that has removed carbon from the atmosphere and combined it with other elements to build their tissue, this vegetation gets compressed by overlying sediments and undergoes chemical transformation.
Deposits of chert formed from the accumulation of siliceous skeletons of microscopic organisms such as radiolaria and diatoms.
Chemical sedimentary rocks
Chemical sedimentary rock forms when mineral constituents in solution become supersaturated and inorganically precipitate. Common chemical sedimentary rocks include oolitic limestone and rocks composed of evaporite minerals, such as halite (rock salt), sylvite, baryte and gypsum.
Other sedimentary rocks
This fourth miscellaneous category includes volcanic tuff and volcanic breccias formed by deposition and later cementation of lava fragments erupted by volcanoes, and impact breccias formed after impact events.
Classification based on composition
Alternatively, sedimentary rocks can be subdivided into compositional groups based on their mineralogy:
Siliciclastic sedimentary rocks, are dominantly composed of silicate minerals. The sediment that makes up these rocks was transported as bed load, suspended load, or by sediment gravity flows. Siliciclastic sedimentary rocks are subdivided into conglomerates and breccias, sandstone, and mudrocks.
Carbonate sedimentary rocks are composed of calcite (rhombohedral ), aragonite (orthorhombic ), dolomite (), and other carbonate minerals based on the ion. Common examples include limestone and the rock dolomite.
Evaporite sedimentary rocks are composed of minerals formed from the evaporation of water. The most common evaporite minerals are carbonates (calcite and others based on ), chlorides (halite and others built on ), and sulfates (gypsum and others built on ). Evaporite rocks commonly include abundant halite (rock salt), gypsum, and anhydrite.
Organic-rich sedimentary rocks have significant amounts of organic material, generally in excess of 3% total organic carbon. Common examples include coal, oil shale as well as source rocks for oil and natural gas.
Siliceous sedimentary rocks are almost entirely composed of silica (), typically as chert, opal, chalcedony or other microcrystalline forms.
Iron-rich sedimentary rocks are composed of >15% iron; the most common forms are banded iron formations and ironstones.
Phosphatic sedimentary rocks are composed of phosphate minerals and contain more than 6.5% phosphorus; examples include deposits of phosphate nodules, bone beds, and phosphatic mudrocks.
Deposition and transformation
Sediment transport and deposition
Sedimentary rocks are formed when sediment is deposited out of air, ice, wind, gravity, or water flows carrying the particles in suspension. This sediment is often formed when weathering and erosion break down a rock into loose material in a source area. The material is then transported from the source area to the deposition area. The type of sediment transported depends on the geology of the hinterland (the source area of the sediment). However, some sedimentary rocks, such as evaporites, are composed of material that form at the place of deposition. The nature of a sedimentary rock, therefore, not only depends on the sediment supply, but also on the sedimentary depositional environment in which it formed.
Transformation (Diagenesis)
As sediments accumulate in a depositional environment, older sediments are buried by younger sediments, and they undergo diagenesis. Diagenesis includes all the chemical, physical, and biological changes, exclusive of surface weathering, undergone by a sediment after its initial deposition. This includes compaction and lithification of the sediments. Early stages of diagenesis, described as eogenesis, take place at shallow depths (a few tens of meters) and is characterized by bioturbation and mineralogical changes in the sediments, with only slight compaction. The red hematite that gives red bed sandstones their color is likely formed during eogenesis. Some biochemical processes, like the activity of bacteria, can affect minerals in a rock and are therefore seen as part of diagenesis.
Deeper burial is accompanied by mesogenesis, during which most of the compaction and lithification takes place. Compaction takes place as the sediments come under increasing overburden (lithostatic) pressure from overlying sediments. Sediment grains move into more compact arrangements, grains of ductile minerals (such as mica) are deformed, and pore space is reduced. Sediments are typically saturated with groundwater or seawater when originally deposited, and as pore space is reduced, much of these connate fluids are expelled. In addition to this physical compaction, chemical compaction may take place via pressure solution. Points of contact between grains are under the greatest strain, and the strained mineral is more soluble than the rest of the grain. As a result, the contact points are dissolved away, allowing the grains to come into closer contact. The increased pressure and temperature stimulate further chemical reactions, such as the reactions by which organic material becomes lignite or coal.
Lithification follows closely on compaction, as increased temperatures at depth hasten the precipitation of cement that binds the grains together. Pressure solution contributes to this process of cementation, as the mineral dissolved from strained contact points is redeposited in the unstrained pore spaces. This further reduces porosity and makes the rock more compact and competent.
Unroofing of buried sedimentary rock is accompanied by telogenesis, the third and final stage of diagenesis. As erosion reduces the depth of burial, renewed exposure to meteoric water produces additional changes to the sedimentary rock, such as leaching of some of the cement to produce secondary porosity.
At sufficiently high temperature and pressure, the realm of diagenesis makes way for metamorphism, the process that forms metamorphic rock.
Properties
Color
The color of a sedimentary rock is often mostly determined by iron, an element with two major oxides: iron(II) oxide and iron(III) oxide. Iron(II) oxide (FeO) only forms under low oxygen (anoxic) circumstances and gives the rock a grey or greenish colour. Iron(III) oxide (Fe2O3) in a richer oxygen environment is often found in the form of the mineral hematite and gives the rock a reddish to brownish colour. In arid continental climates rocks are in direct contact with the atmosphere, and oxidation is an important process, giving the rock a red or orange colour. Thick sequences of red sedimentary rocks formed in arid climates are called red beds. However, a red colour does not necessarily mean the rock formed in a continental environment or arid climate.
The presence of organic material can colour a rock black or grey. Organic material is formed from dead organisms, mostly plants. Normally, such material eventually decays by oxidation or bacterial activity. Under anoxic circumstances, however, organic material cannot decay and leaves a dark sediment, rich in organic material. This can, for example, occur at the bottom of deep seas and lakes. There is little water mixing in such environments; as a result, oxygen from surface water is not brought down, and the deposited sediment is normally a fine dark clay. Dark rocks, rich in organic material, are therefore often shales.
Texture
The size, form and orientation of clasts (the original pieces of rock) in a sediment is called its texture. The texture is a small-scale property of a rock, but determines many of its large-scale properties, such as the density, porosity or permeability.
The 3D orientation of the clasts is called the fabric of the rock. The size and form of clasts can be used to determine the velocity and direction of current in the sedimentary environment that moved the clasts from their origin; fine, calcareous mud only settles in quiet water while gravel and larger clasts are moved only by rapidly moving water. The grain size of a rock is usually expressed with the Wentworth scale, though alternative scales are sometimes used. The grain size can be expressed as a diameter or a volume, and is always an average value, since a rock is composed of clasts with different sizes. The statistical distribution of grain sizes is different for different rock types and is described in a property called the sorting of the rock. When all clasts are more or less of the same size, the rock is called 'well-sorted', and when there is a large spread in grain size, the rock is called 'poorly sorted'.
The form of the clasts can reflect the origin of the rock. For example, coquina, a rock composed of clasts of broken shells, can only form in energetic water. The form of a clast can be described by using four parameters:
Surface texture describes the amount of small-scale relief of the surface of a grain that is too small to influence the general shape. For example, frosted grains, which are covered with small-scale fractures, are characteristic of eolian sandstones.
Rounding describes the general smoothness of the shape of a grain.
Sphericity describes the degree to which the grain approaches a sphere.
Grain form describes the three-dimensional shape of the grain.
Chemical sedimentary rocks have a non-clastic texture, consisting entirely of crystals. To describe such a texture, only the average size of the crystals and the fabric are necessary.
Mineralogy
Most sedimentary rocks contain either quartz (siliciclastic rocks) or calcite (carbonate rocks). In contrast to igneous and metamorphic rocks, a sedimentary rock usually contains very few different major minerals. However, the origin of the minerals in a sedimentary rock is often more complex than in an igneous rock. Minerals in a sedimentary rock may have been present in the original sediments or may formed by precipitation during diagenesis. In the second case, a mineral precipitate may have grown over an older generation of cement. A complex diagenetic history can be established by optical mineralogy, using a petrographic microscope.
Carbonate rocks predominantly consist of carbonate minerals such as calcite, aragonite or dolomite. Both the cement and the clasts (including fossils and ooids) of a carbonate sedimentary rock usually consist of carbonate minerals. The mineralogy of a clastic rock is determined by the material supplied by the source area, the manner of its transport to the place of deposition and the stability of that particular mineral.
The resistance of rock-forming minerals to weathering is expressed by the Goldich dissolution series. In this series, quartz is the most stable, followed by feldspar, micas, and finally other less stable minerals that are only present when little weathering has occurred. The amount of weathering depends mainly on the distance to the source area, the local climate and the time it took for the sediment to be transported to the point where it is deposited. In most sedimentary rocks, mica, feldspar and less stable minerals have been weathered to clay minerals like kaolinite, illite or smectite.
Fossils
Among the three major types of rock, fossils are most commonly found in sedimentary rock. Unlike most igneous and metamorphic rocks, sedimentary rocks form at temperatures and pressures that do not destroy fossil remnants. Often these fossils may only be visible under magnification.
Dead organisms in nature are usually quickly removed by scavengers, bacteria, rotting and erosion, but under exceptional circumstances, these natural processes are unable to take place, leading to fossilisation. The chance of fossilisation is higher when the sedimentation rate is high (so that a carcass is quickly buried), in anoxic environments (where little bacterial activity occurs) or when the organism had a particularly hard skeleton. Larger, well-preserved fossils are relatively rare.
Fossils can be both the direct remains or imprints of organisms and their skeletons. Most commonly preserved are the harder parts of organisms such as bones, shells, and the woody tissue of plants. Soft tissue has a much smaller chance of being fossilized, and the preservation of soft tissue of animals older than 40 million years is very rare. Imprints of organisms made while they were still alive are called trace fossils, examples of which are burrows, footprints, etc.
As a part of a sedimentary rock, fossils undergo the same diagenetic processes as does the host rock. For example, a shell consisting of calcite can dissolve while a cement of silica then fills the cavity. In the same way, precipitating minerals can fill cavities formerly occupied by blood vessels, vascular tissue or other soft tissues. This preserves the form of the organism but changes the chemical composition, a process called permineralization. The most common minerals involved in permineralization are various forms of amorphous silica (chalcedony, flint, chert), carbonates (especially calcite), and pyrite.
At high pressure and temperature, the organic material of a dead organism undergoes chemical reactions in which volatiles such as water and carbon dioxide are expulsed. The fossil, in the end, consists of a thin layer of pure carbon or its mineralized form, graphite. This form of fossilisation is called carbonisation. It is particularly important for plant fossils. The same process is responsible for the formation of fossil fuels like lignite or coal.
Primary sedimentary structures
Structures in sedimentary rocks can be divided into primary structures (formed during deposition) and secondary structures (formed after deposition). Unlike textures, structures are always large-scale features that can easily be studied in the field. Sedimentary structures can indicate something about the sedimentary environment or can serve to tell which side originally faced up where tectonics have tilted or overturned sedimentary layers.
Sedimentary rocks are laid down in layers called beds or strata. A bed is defined as a layer of rock that has a uniform lithology and texture. Beds form by the deposition of layers of sediment on top of each other. The sequence of beds that characterizes sedimentary rocks is called bedding. Single beds can be a couple of centimetres to several meters thick. Finer, less pronounced layers are called laminae, and the structure a lamina forms in a rock is called lamination. Laminae are usually less than a few centimetres thick. Though bedding and lamination are often originally horizontal in nature, this is not always the case. In some environments, beds are deposited at a (usually small) angle. Sometimes multiple sets of layers with different orientations exist in the same rock, a structure called cross-bedding. Cross-bedding is characteristic of deposition by a flowing medium (wind or water).
The opposite of cross-bedding is parallel lamination, where all sedimentary layering is parallel. Differences in laminations are generally caused by cyclic changes in the sediment supply, caused, for example, by seasonal changes in rainfall, temperature or biochemical activity. Laminae that represent seasonal changes (similar to tree rings) are called varves. Any sedimentary rock composed of millimeter or finer scale layers can be named with the general term laminite. When sedimentary rocks have no lamination at all, their structural character is called massive bedding.
Graded bedding is a structure where beds with a smaller grain size occur on top of beds with larger grains. This structure forms when fast flowing water stops flowing. Larger, heavier clasts in suspension settle first, then smaller clasts. Although graded bedding can form in many different environments, it is a characteristic of turbidity currents.
The surface of a particular bed, called the bedform, can also be indicative of a particular sedimentary environment. Examples of bed forms include dunes and ripple marks. Sole markings, such as tool marks and flute casts, are grooves eroded on a surface that are preserved by renewed sedimentation. These are often elongated structures and can be used to establish the direction of the flow during deposition.
Ripple marks also form in flowing water. There can be symmetric or asymmetric. Asymmetric ripples form in environments where the current is in one direction, such as rivers. The longer flank of such ripples is on the upstream side of the current. Symmetric wave ripples occur in environments where currents reverse directions, such as tidal flats.
Mudcracks are a bed form caused by the dehydration of sediment that occasionally comes above the water surface. Such structures are commonly found at tidal flats or point bars along rivers.
Secondary sedimentary structures
Secondary sedimentary structures are those which formed after deposition. Such structures form by chemical, physical and biological processes within the sediment. They can be indicators of circumstances after deposition. Some can be used as way up criteria.
Organic materials in a sediment can leave more traces than just fossils. Preserved tracks and burrows are examples of trace fossils (also called ichnofossils). Such traces are relatively rare. Most trace fossils are burrows of molluscs or arthropods. This burrowing is called bioturbation by sedimentologists. It can be a valuable indicator of the biological and ecological environment that existed after the sediment was deposited. On the other hand, the burrowing activity of organisms can destroy other (primary) structures in the sediment, making a reconstruction more difficult.
Secondary structures can also form by diagenesis or the formation of a soil (pedogenesis) when a sediment is exposed above the water level. An example of a diagenetic structure common in carbonate rocks is a stylolite. Stylolites are irregular planes where material was dissolved into the pore fluids in the rock. This can result in the precipitation of a certain chemical species producing colouring and staining of the rock, or the formation of concretions. Concretions are roughly concentric bodies with a different composition from the host rock. Their formation can be the result of localized precipitation due to small differences in composition or porosity of the host rock, such as around fossils, inside burrows or around plant roots. In carbonate rocks such as limestone or chalk, chert or flint concretions are common, while terrestrial sandstones sometimes contain iron concretions. Calcite concretions in clay containing angular cavities or cracks are called septarian concretions.
After deposition, physical processes can deform the sediment, producing a third class of secondary structures. Density contrasts between different sedimentary layers, such as between sand and clay, can result in flame structures or load casts, formed by inverted diapirism. While the clastic bed is still fluid, diapirism can cause a denser upper layer to sink into a lower layer. Sometimes, density contrasts occur or are enhanced when one of the lithologies dehydrates. Clay can be easily compressed as a result of dehydration, while sand retains the same volume and becomes relatively less dense. On the other hand, when the pore fluid pressure in a sand layer surpasses a critical point, the sand can break through overlying clay layers and flow through, forming discordant bodies of sedimentary rock called sedimentary dykes. The same process can form mud volcanoes on the surface where they broke through upper layers.
Sedimentary dykes can also be formed in a cold climate where the soil is permanently frozen during a large part of the year. Frost weathering can form cracks in the soil that fill with rubble from above. Such structures can be used as climate indicators as well as way up structures.
Density contrasts can also cause small-scale faulting, even while sedimentation progresses (synchronous-sedimentary faulting). Such faulting can also occur when large masses of non-lithified sediment are deposited on a slope, such as at the front side of a delta or the continental slope. Instabilities in such sediments can result in the deposited material to slump, producing fissures and folding. The resulting structures in the rock are syn-sedimentary folds and faults, which can be difficult to distinguish from folds and faults formed by tectonic forces acting on lithified rocks.
Depositional environments
The setting in which a sedimentary rock forms is called the depositional environment. Every environment has a characteristic combination of geologic processes, and circumstances. The type of sediment that is deposited is not only dependent on the sediment that is transported to a place (provenance), but also on the environment itself.
A marine environment means that the rock was formed in a sea or ocean. Often, a distinction is made between deep and shallow marine environments. Deep marine usually refers to environments more than 200 m below the water surface (including the abyssal plain). Shallow marine environments exist adjacent to coastlines and can extend to the boundaries of the continental shelf. The water movements in such environments have a generally higher energy than that in deep environments, as wave activity diminishes with depth. This means that coarser sediment particles can be transported and the deposited sediment can be coarser than in deeper environments. When the sediment is transported from the continent, an alternation of sand, clay and silt is deposited. When the continent is far away, the amount of such sediment deposited may be small, and biochemical processes dominate the type of rock that forms. Especially in warm climates, shallow marine environments far offshore mainly see deposition of carbonate rocks. The shallow, warm water is an ideal habitat for many small organisms that build carbonate skeletons. When these organisms die, their skeletons sink to the bottom, forming a thick layer of calcareous mud that may lithify into limestone. Warm shallow marine environments also are ideal environments for coral reefs, where the sediment consists mainly of the calcareous skeletons of larger organisms.
In deep marine environments, the water current working the sea bottom is small. Only fine particles can be transported to such places. Typically sediments depositing on the ocean floor are fine clay or small skeletons of micro-organisms. At 4 km depth, the solubility of carbonates increases dramatically (the depth zone where this happens is called the lysocline). Calcareous sediment that sinks below the lysocline dissolves; as a result, no limestone can be formed below this depth. Skeletons of micro-organisms formed of silica (such as radiolarians) are not as soluble and are still deposited. An example of a rock formed of silica skeletons is radiolarite. When the bottom of the sea has a small inclination, for example, at the continental slopes, the sedimentary cover can become unstable, causing turbidity currents. Turbidity currents are sudden disturbances of the normally quiet deep marine environment and can cause the near-instantaneous deposition of large amounts of sediment, such as sand and silt. The rock sequence formed by a turbidity current is called a turbidite.
The coast is an environment dominated by wave action. At a beach, dominantly denser sediment such as sand or gravel, often mingled with shell fragments, is deposited, while the silt and clay sized material is kept in mechanical suspension. Tidal flats and shoals are places that sometimes dry because of the tide. They are often cross-cut by gullies, where the current is strong and the grain size of the deposited sediment is larger. Where rivers enter the body of water, either on a sea or lake coast, deltas can form. These are large accumulations of sediment transported from the continent to places in front of the mouth of the river. Deltas are dominantly composed of clastic (rather than chemical) sediment.
A continental sedimentary environment is an environment in the interior of a continent. Examples of continental environments are lagoons, lakes, swamps, floodplains and alluvial fans. In the quiet water of swamps, lakes and lagoons, fine sediment is deposited, mingled with organic material from dead plants and animals. In rivers, the energy of the water is much greater and can transport heavier clastic material. Besides transport by water, sediment can be transported by wind or glaciers. Sediment transported by wind is called aeolian and is almost always very well sorted, while sediment transported by a glacier is called glacial till and is characterized by very poor sorting.
Aeolian deposits can be quite striking. The depositional environment of the Touchet Formation, located in the Northwestern United States, had intervening periods of aridity which resulted in a series of rhythmite layers. Erosional cracks were later infilled with layers of soil material, especially from aeolian processes. The infilled sections formed vertical inclusions in the horizontally deposited layers, and thus provided evidence of the sequence of events during deposition of the forty-one layers of the formation.
Sedimentary facies
The kind of rock formed in a particular depositional environment is called its sedimentary facies. Sedimentary environments usually exist alongside each other in certain natural successions. A beach, where sand and gravel is deposited, is usually bounded by a deeper marine environment a little offshore, where finer sediments are deposited at the same time. Behind the beach, there can be dunes (where the dominant deposition is well sorted sand) or a lagoon (where fine clay and organic material is deposited). Every sedimentary environment has its own characteristic deposits. When sedimentary strata accumulate through time, the environment can shift, forming a change in facies in the subsurface at one location. On the other hand, when a rock layer with a certain age is followed laterally, the lithology (the type of rock) and facies eventually change.
Facies can be distinguished in a number of ways: the most common are by the lithology (for example: limestone, siltstone or sandstone) or by fossil content. Coral, for example, only lives in warm and shallow marine environments and fossils of coral are thus typical for shallow marine facies. Facies determined by lithology are called lithofacies; facies determined by fossils are biofacies.
Sedimentary environments can shift their geographical positions through time. Coastlines can shift in the direction of the sea when the sea level drops (regression), when the surface rises (transgression) due to tectonic forces in the Earth's crust or when a river forms a large delta. In the subsurface, such geographic shifts of sedimentary environments of the past are recorded in shifts in sedimentary facies. This means that sedimentary facies can change either parallel or perpendicular to an imaginary layer of rock with a fixed age, a phenomenon described by Walther's Law.
The situation in which coastlines move in the direction of the continent is called transgression. In the case of transgression, deeper marine facies are deposited over shallower facies, a succession called onlap. Regression is the situation in which a coastline moves in the direction of the sea. With regression, shallower facies are deposited on top of deeper facies, a situation called offlap.
The facies of all rocks of a certain age can be plotted on a map to give an overview of the palaeogeography. A sequence of maps for different ages can give an insight in the development of the regional geography.
Gallery of sedimentary facies
Sedimentary basins
Places where large-scale sedimentation takes place are called sedimentary basins. The amount of sediment that can be deposited in a basin depends on the depth of the basin, the so-called accommodation space. The depth, shape and size of a basin depend on tectonics, movements within the Earth's lithosphere. Where the lithosphere moves upward (tectonic uplift), land eventually rises above sea level and the area becomes a source for new sediment as erosion removes material. Where the lithosphere moves downward (tectonic subsidence), a basin forms and sediments are deposited.
A type of basin formed by the moving apart of two pieces of a continent is called a rift basin. Rift basins are elongated, narrow and deep basins. Due to divergent movement, the lithosphere is stretched and thinned, so that the hot asthenosphere rises and heats the overlying rift basin. Apart from continental sediments, rift basins normally also have part of their infill consisting of volcanic deposits. When the basin grows due to continued stretching of the lithosphere, the rift grows and the sea can enter, forming marine deposits.
When a piece of lithosphere that was heated and stretched cools again, its density rises, causing isostatic subsidence. If this subsidence continues long enough, the basin is called a sag basin. Examples of sag basins are the regions along passive continental margins, but sag basins can also be found in the interior of continents. In sag basins, the extra weight of the newly deposited sediments is enough to keep the subsidence going in a vicious circle. The total thickness of the sedimentary infill in a sag basin can thus exceed 10 km.
A third type of basin exists along convergent plate boundaries – places where one tectonic plate moves under another into the asthenosphere. The subducting plate bends and forms a fore-arc basin in front of the overriding plate – an elongated, deep asymmetric basin. Fore-arc basins are filled with deep marine deposits and thick sequences of turbidites. Such infill is called flysch. When the convergent movement of the two plates results in continental collision, the basin becomes shallower and develops into a foreland basin. At the same time, tectonic uplift forms a mountain belt in the overriding plate, from which large amounts of material are eroded and transported to the basin. Such erosional material of a growing mountain chain is called molasse and has either a shallow marine or a continental facies.
At the same time, the growing weight of the mountain belt can cause isostatic subsidence in the area of the overriding plate on the other side to the mountain belt. The basin type resulting from this subsidence is called a back-arc basin and is usually filled by shallow marine deposits and molasse.
Influence of astronomical cycles
In many cases facies changes and other lithological features in sequences of sedimentary rock have a cyclic nature. This cyclic nature was caused by cyclic changes in sediment supply and the sedimentary environment. Most of these cyclic changes are caused by astronomic cycles. Short astronomic cycles can be the difference between the tides or the spring tide every two weeks. On a larger time-scale, cyclic changes in climate and sea level are caused by Milankovitch cycles: cyclic changes in the orientation and/or position of the Earth's rotational axis and orbit around the Sun. There are a number of Milankovitch cycles known, lasting between 10,000 and 200,000 years.
Relatively small changes in the orientation of the Earth's axis or length of the seasons can be a major influence on the Earth's climate. An example are the ice ages of the past 2.6 million years (the Quaternary period), which are assumed to have been caused by astronomic cycles. Climate change can influence the global sea level (and thus the amount of accommodation space in sedimentary basins) and sediment supply from a certain region. Eventually, small changes in astronomic parameters can cause large changes in sedimentary environment and sedimentation.
Sedimentation rates
The rate at which sediment is deposited differs depending on the location. A channel in a tidal flat can see the deposition of a few metres of sediment in one day, while on the deep ocean floor each year only a few millimetres of sediment accumulate. A distinction can be made between normal sedimentation and sedimentation caused by catastrophic processes. The latter category includes all kinds of sudden exceptional processes like mass movements, rock slides or flooding. Catastrophic processes can see the sudden deposition of a large amount of sediment at once. In some sedimentary environments, most of the total column of sedimentary rock was formed by catastrophic processes, even though the environment is usually a quiet place. Other sedimentary environments are dominated by normal, ongoing sedimentation.
In many cases, sedimentation occurs slowly. In a desert, for example, the wind deposits siliciclastic material (sand or silt) in some spots, or catastrophic flooding of a wadi may cause sudden deposits of large quantities of detrital material, but in most places eolian erosion dominates. The amount of sedimentary rock that forms is not only dependent on the amount of supplied material, but also on how well the material consolidates. Erosion removes most deposited sediment shortly after deposition.
Stratigraphy
Sedimentary rock are laid down in layers called beds or strata, each layer is horizontally laid down over the older ones and new layers are above older layers as stated in the principle of superposition. There are usually some gaps in the sequence called unconformities which represent periods where no new sediments were laid down, or when earlier sedimentary layers were raised above sea level and eroded away.
Unconformities can be classified based on the orientation of the strata on either sides of the unconformity:
Angular unconformity when the earlier layers are tilted and eroded while the later layers are horizontally laid.
Nonconformity if the early layers have no bedding in contrast to the later layers, ie. they are igneous or metamorphic rocks.
Disconformity if both the early beds and the later beds are parallel to each other.
Sedimentary rocks contain important information about the history of the Earth. They contain fossils, the preserved remains of ancient plants and animals. Coal is considered a type of sedimentary rock. The composition of sediments provides us with clues as to the original rock. Differences between successive layers indicate changes to the environment over time. Sedimentary rocks can contain fossils because, unlike most igneous and metamorphic rocks, they form at temperatures and pressures that do not destroy fossil remains.
Provenance
Provenance is the reconstruction of the origin of sediments. All rock exposed at Earth's surface is subjected to physical or chemical weathering and broken down into finer grained sediment. All three types of rocks (igneous, sedimentary and metamorphic rocks) can be the source of sedimentary detritus. The purpose of sedimentary provenance studies is to reconstruct and interpret the history of sediment from the initial parent rocks at a source area to final detritus at a burial place.
| Physical sciences | Petrology | null |
44424 | https://en.wikipedia.org/wiki/Metamorphic%20rock | Metamorphic rock | Metamorphic rocks arise from the transformation of existing rock to new types of rock in a process called metamorphism. The original rock (protolith) is subjected to temperatures greater than and, often, elevated pressure of or more, causing profound physical or chemical changes. During this process, the rock remains mostly in the solid state, but gradually recrystallizes to a new texture or mineral composition. The protolith may be an igneous, sedimentary, or existing metamorphic rock.
Metamorphic rocks make up a large part of the Earth's crust and form 12% of the Earth's land surface. They are classified by their protolith, their chemical and mineral makeup, and their texture. They may be formed simply by being deeply buried beneath the Earth's surface, where they are subject to high temperatures and the great pressure of the rock layers above. They can also form from tectonic processes such as continental collisions, which cause horizontal pressure, friction, and distortion. Metamorphic rock can be formed locally when rock is heated by the intrusion of hot molten rock called magma from the Earth's interior. The study of metamorphic rocks (now exposed at the Earth's surface following erosion and uplift) provides information about the temperatures and pressures that occur at great depths within the Earth's crust.
Some examples of metamorphic rocks are gneiss, slate, marble, schist, and quartzite. Slate and quartzite tiles are used in building construction. Marble is also prized for building construction and as a medium for sculpture. On the other hand, schist bedrock can pose a challenge for civil engineering because of its pronounced planes of weakness.
Origin
Metamorphic rocks form one of the three great divisions of rock types. They are distinguished from igneous rocks, which form from molten magma, and sedimentary rocks, which form from sediments eroded from existing rock or precipitated chemically from bodies of water.
Metamorphic rocks are formed when existing rock is transformed physically or chemically at elevated temperature, without actually melting to any great degree. The importance of heating in the formation of metamorphic rock was first noted by the pioneering Scottish naturalist, James Hutton, who is often described as the father of modern geology. Hutton wrote in 1795 that some rock beds of the Scottish Highlands had originally been sedimentary rock but had been transformed by great heat.
Hutton also speculated that pressure was important in metamorphism. This hypothesis was tested by his friend, James Hall, who sealed chalk into a makeshift pressure vessel constructed from a cannon barrel and heated it in an iron foundry furnace. Hall found that this produced a material strongly resembling marble, rather than the usual quicklime produced by heating of chalk in the open air. French geologists subsequently added metasomatism, the circulation of fluids through buried rock, to the list of processes that help bring about metamorphism. However, metamorphism can take place without metasomatism (isochemical metamorphism) or at depths of just a few hundred meters where pressures are relatively low (for example, in contact metamorphism).
Metamorphic processes change the texture or mineral composition of the metamorphosed rock.
Mineralogical changes
Metasomatism can change the bulk composition of a rock. Hot fluids circulating through pore space in the rock can dissolve existing minerals and precipitate new minerals. Dissolved substances are transported out of the rock by the fluids while new substances are brought in by fresh fluids. This can obviously change the mineral makeup of the rock.
However, changes in the mineral composition can take place even when the bulk composition of the rock does not change. This is possible because all minerals are stable only within certain limits of temperature, pressure, and chemical environment. For example, at atmospheric pressure, the mineral kyanite transforms to andalusite at a temperature of about . Andalusite, in turn, transforms to sillimanite when the temperature reaches about . All three have the identical composition, . Likewise, forsterite is stable over a broad range of pressure and temperature in marble, but is converted to pyroxene at elevated pressure and temperature in more silicate-rich rock containing plagioclase, with which the forsterite reacts chemically.
Many complex high-temperature reactions may take place between minerals without them melting, and each mineral assemblage produced indicates the temperatures and pressures at the time of metamorphism. These reactions are possible because of rapid diffusion of atoms at elevated temperature. Pore fluid between mineral grains can be an important medium through which atoms are exchanged.
Textural changes
The change in the particle size of the rock during the process of metamorphism is called recrystallization. For instance, the small calcite crystals in the sedimentary rock limestone and chalk change into larger crystals in the metamorphic rock marble. In metamorphosed sandstone, recrystallization of the original quartz sand grains results in very compact quartzite, also known as metaquartzite, in which the often larger quartz crystals are interlocked. Both high temperatures and pressures contribute to recrystallization. High temperatures allow the atoms and ions in solid crystals to migrate, thus reorganizing the crystals, while high pressures cause solution of the crystals within the rock at their point of contact.
Description
Metamorphic rocks are characterized by their distinctive mineral composition and texture.
Metamorphic minerals
Because every mineral is stable only within certain limits, the presence of certain minerals in metamorphic rocks indicates the approximate temperatures and pressures at which the rock underwent metamorphism. These minerals are known as index minerals. Examples include sillimanite, kyanite, staurolite, andalusite, and some garnet.
Other minerals, such as olivines, pyroxenes, hornblende, micas, feldspars, and quartz, may be found in metamorphic rocks but are not necessarily the result of the process of metamorphism. These minerals can also form during the crystallization of igneous rocks. They are stable at high temperatures and pressures and may remain chemically unchanged during the metamorphic process.
Texture
Metamorphic rocks are typically more coarsely crystalline than the protolith from which they formed. Atoms in the interior of a crystal are surrounded by a stable arrangement of neighboring atoms. This is partially missing at the surface of the crystal, producing a surface energy that makes the surface thermodynamically unstable. Recrystallization to coarser crystals reduces the surface area and so minimizes the surface energy.
Although grain coarsening is a common result of metamorphism, rock that is intensely deformed may eliminate strain energy by recrystallizing as a fine-grained rock called mylonite. Certain kinds of rock, such as those rich in quartz, carbonate minerals, or olivine, are particularly prone to form mylonites, while feldspar and garnet are resistant to mylonitization.
Foliation
Many kinds of metamorphic rocks show a distinctive layering called foliation (derived from the Latin word folia, meaning "leaves"). Foliation develops when a rock is being shortened along one axis during recrystallization. This causes crystals of platy minerals, such as mica and chlorite, to become rotated such that their short axes are parallel to the direction of shortening. This results in a banded, or foliated, rock, with the bands showing the colors of the minerals that formed them. Foliated rock often develops planes of cleavage. Slate is an example of a foliated metamorphic rock, originating from shale, and it typically shows well-developed cleavage that allows slate to be split into thin plates.
The type of foliation that develops depends on the metamorphic grade. For instance, starting with a mudstone, the following sequence develops with increasing temperature: The mudstone is first converted to slate, which is a very fine-grained, foliated metamorphic rock, characteristic of very low grade metamorphism. Slate in turn is converted to phyllite, which is fine-grained and found in areas of low grade metamorphism. Schist is medium to coarse-grained and found in areas of medium grade metamorphism. High-grade metamorphism transforms the rock to gneiss, which is coarse to very coarse-grained.
Rocks that were subjected to uniform pressure from all sides, or those that lack minerals with distinctive growth habits, will not be foliated. Marble lacks platy minerals and is generally not foliated, which allows its use as a material for sculpture and architecture.
Classification
Metamorphic rocks are one of the three great divisions of all rock types, and so there is a great variety of metamorphic rock types. In general, if the protolith of a metamorphic rock can be determined, the rock is described by adding the prefix meta- to the protolith rock name. For example, if the protolith is known to be basalt, the rock will be described as a metabasalt. Likewise, a metamorphic rock whose protolith is known to be a conglomerate will be described as a metaconglomerate. For a metamorphic rock to be classified in this manner, the protolith should be identifiable from the characteristics of the metamorphic rock itself, and not inferred from other information.
Under the British Geological Survey's classification system, if all that can be determined about the protolith is its general type, such as sedimentary or volcanic, the classification is based on the mineral mode (the volume percentages of different minerals in the rock). Metasedimentary rocks are divided into carbonate-rich rock (metacarbonates or calcsilicate-rocks) or carbonate-poor rocks, and the latter are further classified by the relative abundance of mica in their composition. This ranges from low-mica psammite through semipelite to high-mica pelite. Psammites composed mostly of quartz are classified as quartzite. Metaigneous rocks are classified similarly to igneous rocks, by silica content, from meta-ultramafic-rock (which is very low in silica) to metafelsic-rock (with a high silica content).
Where the mineral mode cannot be determined, as is often the case when rock is first examined in the field, then classification must be based on texture. The textural types are:
Schists, which are medium-grained strongly foliated rocks. These show the most well-developed schistosity, defined as the extent to which platy minerals are present and are aligned in a single direction, so that the rock easily splits into plates less than a centimeter (0.4 inches) thick.
Gneisses, which are more coarse grained and show thicker foliation than schists, with layers over 5mm thick. These show less well-developed schistosity.
Granofels, which show no obvious foliation or schistosity.
A hornfels is a granofels that is known to result from contact metamorphism. A slate is a fine-grained metamorphic rock that easily splits into thin plates but shows no obvious compositional layering. The term is used only when very little else is known about the rock that would allow a more definite classification. Textural classifications may be prefixed to indicate a sedimentary protolith (para-, such as paraschist) or igneous protolith (ortho-, such as orthogneiss). When nothing is known about the protolith, the textural name is used without a prefix. For example, a schist is a rock with schistose texture whose protolith is uncertain.
Special classifications exist for metamorphic rocks with a volcaniclastic protolith or formed along a fault or through hydrothermal circulation. A few special names are used for rocks of unknown protolith but known modal composition, such as marble, eclogite, or amphibolite. Special names may also be applied more generally to rocks dominated by a single mineral, or with a distinctive composition or mode or origin. Special names still in wide use include amphibolite, greenschist, phyllite, marble, serpentinite, eclogite, migmatite, skarn, granulite, mylonite, and slate.
The basic classification can be supplemented by terms describing mineral content or texture. For example, a metabasalt showing weak schistosity might be described as a gneissic metabasalt, and a pelite containing abundant staurolite might be described as a staurolite pelite.
Metamorphic facies
A metamorphic facies is a set of distinctive assemblages of minerals that are found in metamorphic rock that formed under a specific combination of pressure and temperature. The particular assemblage is somewhat dependent on the composition of that protolith, so that (for example) the amphibolite facies of a marble will not be identical with the amphibolite facies of a pelite. However, the facies are defined such that metamorphic rock with as broad a range of compositions as is practical can be assigned to a particular facies. The present definition of metamorphic facies is largely based on the work of the Finnish geologist, Pentti Eskola, with refinements based on subsequent experimental work. Eskola drew upon the zonal schemes, based on index minerals, that were pioneered by the British geologist, George Barrow.
The metamorphic facies is not usually considered when classifying metamorphic rock based on protolith, mineral mode, or texture. However, a few metamorphic facies produce rock of such distinctive character that the facies name is used for the rock when more precise classification is not possible. The chief examples are amphibolite and eclogite. The British Geological Survey strongly discourages the use of granulite as a classification for rock metamorphosed to the granulite facies. Instead, such rock will often be classified as a granofels. However, this approach is not universally accepted.
Occurrence
Metamorphic rocks make up a large part of the Earth's crust and form 12% of the Earth's land surface. The lower continental crust is mostly metamafic-rock and pelite which have reached the granulite facies. The middle continental crust is dominated by metamorphic rock that has reached the amphibolite facies. Within the upper crust, which is the only part of the Earth's crust geologists can directly sample, metamorphic rock forms only from processes that can occur at shallow depth. These are contact (thermal) metamorphism, dynamic (cataclastic) metamorphism, hydrothermal metamorphism, and impact metamorphism. These processes are relatively local in occurrence and usually reach only the low-pressure facies, such as the hornfels and sanidinite facies. Most metamorphic rock is formed by regional metamorphism in the middle and lower crust, where the rock reaches the higher-pressure metamorphic facies. This rock is found at the surface only where extensive uplift and erosion has exhumed rock that was formerly much deeper in the crust.
Orogenic belts
Metamorphic rock is extensively exposed in orogenic belts produced by the collision of tectonic plates at convergent boundaries. Here formerly deeply buried rock has been brought to the surface by uplift and erosion. The metamorphic rock exposed in orogenic belts may have been metamorphosed simply by being at great depths below the Earth's surface, subjected to high temperatures and the great pressure caused by the immense weight of the rock layers above. This kind of regional metamorphism is known as burial metamorphism. This tends to produce low-grade metamorphic rock. Much more common is metamorphic rock formed during the collision process itself. The collision of plates causes high temperatures, pressures and deformation in the rocks along these belts. Metamorphic rock formed in these settings tends to shown well-developed schistosity.
Metamorphic rock of orogenic belts shows a variety of metamorphic facies. Where subduction is taking place, the basalt of the subducting slab is metamorphosed to high-pressure metamorphic facies. It initially undergoes low-grade metamorphism to metabasalt of the zeolite and prehnite-pumpellyite facies, but as the basalt subducts to greater depths, it is metamorphosed to the blueschist facies and then the eclogite facies. Metamorphism to the eclogite facies releases a great deal of water vapor from the rock, which drives volcanism in the overlying volcanic arc. Eclogite is also significantly denser than blueschist, which drives further subduction of the slab deep into the Earth's mantle. Metabasalt and blueschist may be preserved in blueschist metamorphic belts formed by collisions between continents. They may also be preserved by obduction onto the overriding plate as part of ophiolites. Eclogites are occasionally found at sites of continental collision, where the subducted rock is rapidly brought back to the surface, before it can be converted to the granulite facies in the hot upper mantle. Many samples of eclogite are xenoliths brought to the surface by volcanic activity.
Many orogenic belts contain higher-temperature, lower-pressure metamorphic belts. These may form through heating of the rock by ascending magmas of volcanic arcs, but on a regional scale. Deformation and crustal thickening in an orogenic belt may also produce these kinds of metamorphic rocks. These rocks reach the greenschist, amphibolite, or granulite facies and are the most common of metamorphic rocks produced by regional metamorphosis. The association of an outer high-pressure, low-temperature metamorphic zone with an inner zone of low-pressure, high-temperature metamorphic rocks is called a paired metamorphic belt. The main islands of Japan show three distinct paired metamorphic belts, corresponding to different episodes of subduction.
Metamorphic core complexes
Metamorphic rock is also exposed in metamorphic core complexes, which form in region of crustal extension. They are characterized by low-angle faulting that exposes domes of middle or lower crust metamorphic rock. These were first recognized and studied in the Basin and Range Province of southwestern North America, but are also found in southern Aegean Sea, in the D'Entrecasteaux Islands, and in other areas of extension.
Granite-greenstone belts
Continental shields are regions of exposed ancient rock that make up the stable cores of continents. The rock exposed in the oldest regions of shields, which is of Archean age (over 2500 million years old), mostly belong to granite-greenstone belts. The greenstone belts contain metavolcanic and metasedimentary rock that has undergone a relatively mild grade of metamorphism, at temperatures of and pressures of . They can be divided into a lower group of metabasalts, including rare metakomatiites; a middle group of meta-intermediate-rock and meta-felsic-rock; and an upper group of metasedimentary rock.
The greenstone belts are surrounded by high-grade gneiss terrains showing highly deformed low-pressure, high-temperature (over ) metamorphism to the amphibolite or granulite facies. These form most of the exposed rock in Archean cratons.
The granite-greenstone belts are intruded by a distinctive group of granitic rocks called the tonalite-trondhjemite-granodiorite or TTG suite. These are the most voluminous rocks in the craton and may represent an important early phase in the formation of continental crust.
Mid-ocean ridges
Mid-ocean ridges are where new oceanic crust is formed as tectonic plates move apart. Hydrothermal metamorphism is extensive here. This is characterized by metasomatism by hot fluids circulating through the rock. This produces metamorphic rock of the greenschist facies. The metamorphic rock, serpentinite, is particularly characteristic of these settings, and represents chemical transformation of olivine and pyroxene in ultramafic rock to serpentine group minerals.
Contact aureoles
Contact metamorphism takes place when magma is injected into the surrounding solid rock (country rock). The changes that occur are greatest wherever the magma comes into contact with the rock because the temperatures are highest at this boundary and decrease with distance from it. Around the igneous rock that forms from the cooling magma is a metamorphosed zone called a contact aureole. Aureoles may show all degrees of metamorphism from the contact area to unmetamorphosed (unchanged) country rock some distance away. The formation of important ore minerals may occur by the process of metasomatism at or near the contact zone. Contact aureoles around large plutons may be as much as several kilometers wide.
The term hornfels is often used by geologists to signify those fine grained, compact, non-foliated products of contact metamorphism. The contact aureole typically shows little deformation, and so hornfels is usually devoid of schistosity and forms a tough, equigranular rock. If the rock was originally banded or foliated (as, for example, a laminated sandstone or a foliated calc-schist) this character may not be obliterated, and a banded hornfels is the product. Contact metamorphism close to the surface produces distinctive low-pressure metamorphic minerals, such as spinel, andalusite, vesuvianite, or wollastonite.
Similar changes may be induced in shales by the burning of coal seams. This produces a rock type named clinker.
There is also a tendency for metasomatism between the igneous magma and sedimentary country rock, whereby the chemicals in each are exchanged or introduced into the other. In that case, hybrid rocks called skarn arise.
Other occurrences
Dynamic (cataclastic) metamorphism takes place locally along faults. Here intense shearing of the rock typically forms mylonites.
Impact metamorphism is unlike other forms of metamorphism in that it takes place during impact events by extraterrestrial bodies. It produces rare ultrahigh pressure metamorphic minerals, such as coesite and stishovite. Coesite is rarely found in eclogite brought to the surface in kimberlite pipes, but the presence of stishovite is unique to impact structures.
Uses
Slate tiles are used in construction, particularly as roof shingle.
Quartzite is sufficiently hard and dense that it is difficult to quarry. However, some quartzite is used as dimension stone, often as slabs for flooring, walls, or stairsteps. About 6% of crushed stone, used mostly for road aggregate, is quartzite.
Marble is also prized for building construction and as a medium for sculpture.
Hazards
Schistose bedrock can pose a challenge for civil engineering because of its pronounced planes of weakness. A hazard may exist even in undisturbed terrain. On August 17, 1959, a magnitude 7.2 earthquake destabilized a mountain slope near Hebgen Lake, Montana, composed of schist. This caused a massive landslide that killed 26 people camping in the area.
Metamorphosed ultramafic rock contains serpentine group minerals, which includes varieties of asbestos that pose a hazard to human health.
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44430 | https://en.wikipedia.org/wiki/Glacis | Glacis | A glacis (, ) in military engineering is an artificial slope as part of a medieval castle or in early modern fortresses. They may be constructed of earth as a temporary structure or of stone in more permanent structure. More generally, a glacis is any slope, natural or artificial, which fulfils the above requirements. The etymology of this French word suggests a slope made dangerous with ice, hence the relationship with glacier.
A glacis plate is the sloped front-most section of the hull of a tank or other armoured fighting vehicle.
Ancient fortifications
A glacis could also appear in ancient fortresses, such as the one the ancient Egyptians built at Semna in Nubia. Here it was used by them to prevent enemy siege engines from weakening defensive walls.
Hillforts in Britain started to incorporate glacis around 350 BC. Those at Maiden Castle, Dorset were high.
Medieval fortifications
Glacises, also called taluses, were incorporated into medieval fortifications to strengthen the walls against undermining, to hamper escalades and so that missiles dropped from the battlements would ricochet off the glacis into attacking forces.
Towards the end of the medieval period some castles were modified to make them defensible against cannons. Glacis consisting of earthen slopes faced with stones were placed in front of the curtain walls and bastions (towers) to absorb the impact of cannon shots or to deflect them. Towers were lowered to the same height as the curtain walls and converted into gun platforms.
Early modern European fortifications
Early modern European fortresses were so constructed as to keep any potential assailant under the fire of the defenders until the last possible moment. On natural, level ground, troops attacking any high work have a degree of shelter from its fire when close up to it; the glacis consists of a slope with a low grade inclined towards the top of the wall. This gave defenders a direct line of sight into the assaulting force, allowing them to efficiently sweep the field with fire from the parapet. Additionally, but secondarily, the bank of earth would shield the walls from being hit directly by cannon fire.
Though defenders on high ground already have a direct line of sight, a glacis allows the field of fire to be swept more efficiently by minimizing changes to the angle of their guns while firing. Furthermore, the glacis prevents attacking cannon from having a clear shot at the walls of a fortress, as usually these cannot be seen until the glacis is crossed and the ditch, bounded on either side by the smooth, masoned scarp and counterscarp, is reached.
Armored vehicles
The term glacis plate describes the sloped front-most section of the hull of a tank or other armored fighting vehicle, often composed of upper and lower halves. In a head-on-head armored engagement, the glacis plate is the largest and most obvious target available to an enemy gunner.
Sloped armour has two advantages: many projectiles will deflect rather than penetrate; those that attempt to will have to travel on a longer diagonal route through any given thickness of armor than if it were perpendicular to their trajectory.
Anti-tank mines that employ a tilt-rod fuze are also designed to detonate directly underneath the glacis plate. As a result, it is generally the thickest, most robust armored section of a tank, followed by the turret face and gun mantlet.
Gallery
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44469 | https://en.wikipedia.org/wiki/Pluto | Pluto | Pluto (minor-planet designation: 134340 Pluto) is a dwarf planet in the Kuiper belt, a ring of bodies beyond the orbit of Neptune. It is the ninth-largest and tenth-most-massive known object to directly orbit the Sun. It is the largest known trans-Neptunian object by volume, by a small margin, but is less massive than Eris. Like other Kuiper belt objects, Pluto is made primarily of ice and rock and is much smaller than the inner planets. Pluto has roughly one-sixth the mass of the Moon, and one-third its volume.
Pluto has a moderately eccentric and inclined orbit, ranging from from the Sun. Light from the Sun takes 5.5 hours to reach Pluto at its orbital distance of . Pluto's eccentric orbit periodically brings it closer to the Sun than Neptune, but a stable orbital resonance prevents them from colliding.
Pluto has five known moons: Charon, the largest, whose diameter is just over half that of Pluto; Styx; Nix; Kerberos; and Hydra. Pluto and Charon are sometimes considered a binary system because the barycenter of their orbits does not lie within either body, and they are tidally locked. New Horizons was the first spacecraft to visit Pluto and its moons, making a flyby on July 14, 2015, and taking detailed measurements and observations.
Pluto was discovered in 1930 by Clyde W. Tombaugh, making it by far the first known object in the Kuiper belt. It was immediately hailed as the ninth planet. However, its planetary status was questioned when it was found to be much smaller than expected. These doubts increased following the discovery of additional objects in the Kuiper belt starting in the 1990s, and particularly the more massive scattered disk object Eris in 2005. In 2006, the International Astronomical Union (IAU) formally redefined the term planet to exclude dwarf planets such as Pluto. Many planetary astronomers, however, continue to consider Pluto and other dwarf planets to be planets.
History
Discovery
In the 1840s, Urbain Le Verrier used Newtonian mechanics to predict the position of the then-undiscovered planet Neptune after analyzing perturbations in the orbit of Uranus. Subsequent observations of Neptune in the late 19th century led astronomers to speculate that Uranus's orbit was being disturbed by another planet besides Neptune.
In 1906, Percival Lowell—a wealthy Bostonian who had founded Lowell Observatory in Flagstaff, Arizona, in 1894—started an extensive project in search of a possible ninth planet, which he termed "Planet X". By 1909, Lowell and William H. Pickering had suggested several possible celestial coordinates for such a planet. Lowell and his observatory conducted his search, using mathematical calculations made by Elizabeth Williams, until his death in 1916, but to no avail. Unknown to Lowell, his surveys had captured two faint images of Pluto on March 19 and April 7, 1915, but they were not recognized for what they were. There are fourteen other known precovery observations, with the earliest made by the Yerkes Observatory on August 20, 1909.
Percival's widow, Constance Lowell, entered into a ten-year legal battle with the Lowell Observatory over her husband's legacy, and the search for Planet X did not resume until 1929. Vesto Melvin Slipher, the observatory director, gave the job of locating Planet X to 23-year-old Clyde Tombaugh, who had just arrived at the observatory after Slipher had been impressed by a sample of his astronomical drawings.
Tombaugh's task was to systematically image the night sky in pairs of photographs, then examine each pair and determine whether any objects had shifted position. Using a blink comparator, he rapidly shifted back and forth between views of each of the plates to create the illusion of movement of any objects that had changed position or appearance between photographs. On February 18, 1930, after nearly a year of searching, Tombaugh discovered a possible moving object on photographic plates taken on January 23 and 29. A lesser-quality photograph taken on January 21 helped confirm the movement. After the observatory obtained further confirmatory photographs, news of the discovery was telegraphed to the Harvard College Observatory on March 13, 1930.
One Plutonian year corresponds to 247.94 Earth years; thus, in 2178, Pluto will complete its first orbit since its discovery.
Name and symbol
The name Pluto came from the Roman god of the underworld; and it is also an epithet for Hades (the Greek equivalent of Pluto).
Upon the announcement of the discovery, Lowell Observatory received over a thousand suggestions for names. Three names topped the list: Minerva, Pluto and Cronus. 'Minerva' was the Lowell staff's first choice but was rejected because it had already been used for an asteroid; Cronus was disfavored because it was promoted by an unpopular and egocentric astronomer, Thomas Jefferson Jackson See. A vote was then taken and 'Pluto' was the unanimous choice. To make sure the name stuck, and that the planet would not suffer changes in its name as Uranus had, Lowell Observatory proposed the name to the American Astronomical Society and the Royal Astronomical Society; both approved it unanimously. The name was published on May 1, 1930.
The name Pluto had received some 150 nominations among the letters and telegrams sent to Lowell. The first had been from Venetia Burney (1918–2009), an eleven-year-old schoolgirl in Oxford, England, who was interested in classical mythology. She had suggested it to her grandfather Falconer Madan when he read the news of Pluto's discovery to his family over breakfast; Madan passed the suggestion to astronomy professor Herbert Hall Turner, who cabled it to colleagues at Lowell on March 16, three days after the announcement.
The name 'Pluto' was mythologically appropriate: the god Pluto was one of six surviving children of Saturn, and the others had already all been chosen as names of major or minor planets (his brothers Jupiter and Neptune, and his sisters Ceres, Juno and Vesta). Both the god and the planet inhabited "gloomy" regions, and the god was able to make himself invisible, as the planet had been for so long.
The choice was further helped by the fact that the first two letters of Pluto were the initials of Percival Lowell; indeed, 'Percival' had been one of the more popular suggestions for a name for the new planet.
Pluto's planetary symbol was then created as a monogram of the letters "PL". This symbol is rarely used in astronomy anymore, though it is still common in astrology. However, the most common astrological symbol for Pluto, occasionally used in astronomy as well, is an orb (possibly representing Pluto's invisibility cap) over Pluto's bident , which dates to the early 1930s.
The name 'Pluto' was soon embraced by wider culture. In 1930, Walt Disney was apparently inspired by it when he introduced for Mickey Mouse a canine companion named Pluto, although Disney animator Ben Sharpsteen could not confirm why the name was given. In 1941, Glenn T. Seaborg named the newly created element plutonium after Pluto, in keeping with the tradition of naming elements after newly discovered planets, following uranium, which was named after Uranus, and neptunium, which was named after Neptune.
Most languages use the name "Pluto" in various transliterations. In Japanese, Houei Nojiri suggested the calque , and this was borrowed into Chinese and Korean. Some languages of India use the name Pluto, but others, such as Hindi, use the name of Yama, the God of Death in Hinduism. Polynesian languages also tend to use the indigenous god of the underworld, as in Māori Whiro.
Vietnamese might be expected to follow Chinese, but does not because the Sino-Vietnamese word 冥 minh "dark" is homophonous with 明 minh "bright". Vietnamese instead uses Yama, which is also a Buddhist deity, in the form of Sao Diêm Vương 星閻王 "Yama's Star", derived from Chinese 閻王 Yán Wáng / Yìhm Wòhng "King Yama".
Planet X disproved
Once Pluto was found, its faintness and lack of a viewable disc cast doubt on the idea that it was Lowell's Planet X. Estimates of Pluto's mass were revised downward throughout the 20th century.
Astronomers initially calculated its mass based on its presumed effect on Neptune and Uranus. In 1931, Pluto was calculated to be roughly the mass of Earth, with further calculations in 1948 bringing the mass down to roughly that of Mars. In 1976, Dale Cruikshank, Carl Pilcher and David Morrison of the University of Hawaiʻi calculated Pluto's albedo for the first time, finding that it matched that for methane ice; this meant Pluto had to be exceptionally luminous for its size and therefore could not be more than 1 percent the mass of Earth. (Pluto's albedo is times that of Earth.)
In 1978, the discovery of Pluto's moon Charon allowed the measurement of Pluto's mass for the first time: roughly 0.2% that of Earth, and far too small to account for the discrepancies in the orbit of Uranus. Subsequent searches for an alternative Planet X, notably by Robert Sutton Harrington, failed. In 1992, Myles Standish used data from Voyager 2'''s flyby of Neptune in 1989, which had revised the estimates of Neptune's mass downward by 0.5%—an amount comparable to the mass of Mars—to recalculate its gravitational effect on Uranus. With the new figures added in, the discrepancies, and with them the need for a Planet X, vanished. the majority of scientists agree that Planet X, as Lowell defined it, does not exist. Lowell had made a prediction of Planet X's orbit and position in 1915 that was fairly close to Pluto's actual orbit and its position at that time; Ernest W. Brown concluded soon after Pluto's discovery that this was a coincidence.
Classification
From 1992 onward, many bodies were discovered orbiting in the same volume as Pluto, showing that Pluto is part of a population of objects called the Kuiper belt. This made its official status as a planet controversial, with many questioning whether Pluto should be considered together with or separately from its surrounding population. Museum and planetarium directors occasionally created controversy by omitting Pluto from planetary models of the Solar System. In February 2000 the Hayden Planetarium in New York City displayed a Solar System model of only eight planets, which made headlines almost a year later.
Ceres, Pallas, Juno and Vesta lost their planet status among most astronomers after the discovery of many other asteroids in the 1840s. On the other hand, planetary geologists often regarded Ceres, and less often Pallas and Vesta, as being different from smaller asteroids because they were large enough to have undergone geological evolution. Although the first Kuiper belt objects discovered were quite small, objects increasingly closer in size to Pluto were soon discovered, some large enough (like Pluto itself) to satisfy geological but not dynamical ideas of planethood. On July 29, 2005, the debate became unavoidable when astronomers at Caltech announced the discovery of a new trans-Neptunian object, Eris, which was substantially more massive than Pluto and the most massive object discovered in the Solar System since Triton in 1846. Its discoverers and the press initially called it the tenth planet, although there was no official consensus at the time on whether to call it a planet. Others in the astronomical community considered the discovery the strongest argument for reclassifying Pluto as a minor planet.
IAU classification
The debate came to a head in August 2006, with an IAU resolution that created an official definition for the term "planet". According to this resolution, there are three conditions for an object in the Solar System to be considered a planet:
The object must be in orbit around the Sun.
The object must be massive enough to be rounded by its own gravity. More specifically, its own gravity should pull it into a shape defined by hydrostatic equilibrium.
It must have cleared the neighborhood around its orbit.
Pluto fails to meet the third condition. Its mass is substantially less than the combined mass of the other objects in its orbit: 0.07 times, in contrast to Earth, which is 1.7 million times the remaining mass in its orbit (excluding the moon). The IAU further decided that bodies that, like Pluto, meet criteria 1 and 2, but do not meet criterion 3 would be called dwarf planets. In September 2006, the IAU included Pluto, and Eris and its moon Dysnomia, in their Minor Planet Catalogue, giving them the official minor-planet designations "(134340) Pluto", "(136199) Eris", and "(136199) Eris I Dysnomia". Had Pluto been included upon its discovery in 1930, it would have likely been designated 1164, following 1163 Saga, which was discovered a month earlier.
There has been some resistance within the astronomical community toward the reclassification, and in particular planetary scientists often continue to reject it, considering Pluto, Charon, and Eris to be planets for the same reason they do so for Ceres. In effect, this amounts to accepting only the second clause of the IAU definition. Alan Stern, principal investigator with NASA's New Horizons mission to Pluto, derided the IAU resolution. He also stated that because less than five percent of astronomers voted for it, the decision was not representative of the entire astronomical community. Marc W. Buie, then at the Lowell Observatory, petitioned against the definition. Others have supported the IAU, for example Mike Brown, the astronomer who discovered Eris.
Public reception to the IAU decision was mixed. A resolution introduced in the California State Assembly facetiously called the IAU decision a "scientific heresy". The New Mexico House of Representatives passed a resolution in honor of Clyde Tombaugh, the discoverer of Pluto and a longtime resident of that state, that declared that Pluto will always be considered a planet while in New Mexican skies and that March 13, 2007, was Pluto Planet Day. The Illinois Senate passed a similar resolution in 2009 on the basis that Tombaugh was born in Illinois. The resolution asserted that Pluto was "unfairly downgraded to a 'dwarf' planet" by the IAU." Some members of the public have also rejected the change, citing the disagreement within the scientific community on the issue, or for sentimental reasons, maintaining that they have always known Pluto as a planet and will continue to do so regardless of the IAU decision. In 2006, in its 17th annual words-of-the-year vote, the American Dialect Society voted plutoed as the word of the year. To "pluto" is to "demote or devalue someone or something". In April 2024, Arizona (where Pluto was first discovered in 1930) passed a law naming Pluto as the official state planet.
Researchers on both sides of the debate gathered in August 2008, at the Johns Hopkins University Applied Physics Laboratory for a conference that included back-to-back talks on the IAU definition of a planet. Entitled "The Great Planet Debate", the conference published a post-conference press release indicating that scientists could not come to a consensus about the definition of planet. In June 2008, the IAU had announced in a press release that the term "plutoid" would henceforth be used to refer to Pluto and other planetary-mass objects that have an orbital semi-major axis greater than that of Neptune, though the term has not seen significant use.
Orbit
Pluto's orbital period is about 248 years. Its orbital characteristics are substantially different from those of the planets, which follow nearly circular orbits around the Sun close to a flat reference plane called the ecliptic. In contrast, Pluto's orbit is moderately inclined relative to the ecliptic (over 17°) and moderately eccentric (elliptical). This eccentricity means a small region of Pluto's orbit lies closer to the Sun than Neptune's. The Pluto–Charon barycenter came to perihelion on September 5, 1989, and was last closer to the Sun than Neptune between February 7, 1979, and February 11, 1999.
Although the 3:2 resonance with Neptune (see below) is maintained, Pluto's inclination and eccentricity behave in a chaotic manner. Computer simulations can be used to predict its position for several million years (both forward and backward in time), but after intervals much longer than the Lyapunov time of 10–20 million years, calculations become unreliable: Pluto is sensitive to immeasurably small details of the Solar System, hard-to-predict factors that will gradually change Pluto's position in its orbit.
The semi-major axis of Pluto's orbit varies between about 39.3 and 39.6 AU with a period of about 19,951 years, corresponding to an orbital period varying between 246 and 249 years. The semi-major axis and period are presently getting longer.
Relationship with Neptune
Despite Pluto's orbit appearing to cross that of Neptune when viewed from north or south of the Solar System, the two objects' orbits do not intersect. When Pluto is closest to the Sun, and close to Neptune's orbit as viewed from such a position, it is also the farthest north of Neptune's path. Pluto's orbit passes about 8 AU north of that of Neptune, preventing a collision.
This alone is not enough to protect Pluto; perturbations from the planets (especially Neptune) could alter Pluto's orbit (such as its orbital precession) over millions of years so that a collision could happen. However, Pluto is also protected by its 2:3 orbital resonance with Neptune: for every two orbits that Pluto makes around the Sun, Neptune makes three, in a frame of reference that rotates at the rate that Pluto's perihelion precesses (about degrees per year). Each cycle lasts about 495 years. (There are many other objects in this same resonance, called plutinos.) At present, in each 495-year cycle, the first time Pluto is at perihelion (such as in 1989), Neptune is 57° ahead of Pluto. By Pluto's second passage through perihelion, Neptune will have completed a further one and a half of its own orbits, and will be 123° behind Pluto. Pluto and Neptune's minimum separation is over 17 AU, which is greater than Pluto's minimum separation from Uranus (11 AU). The minimum separation between Pluto and Neptune actually occurs near the time of Pluto's aphelion.
The 2:3 resonance between the two bodies is highly stable and has been preserved over millions of years. This prevents their orbits from changing relative to one another, so the two bodies can never pass near each other. Even if Pluto's orbit were not inclined, the two bodies could never collide. When Pluto's period is slightly different from 3/2 of Neptune's, the pattern of its distance from Neptune will drift. Near perihelion Pluto moves interior to Neptune's orbit and is therefore moving faster, so during the first of two orbits in the 495-year cycle, it is approaching Neptune from behind. At present it remains between 50° and 65° behind Neptune for 100 years (e.g. 1937–2036). The gravitational pull between the two causes angular momentum to be transferred to Pluto. This situation moves Pluto into a slightly larger orbit, where it has a slightly longer period, according to Kepler's third law. After several such repetitions, Pluto is sufficiently delayed that at the second perihelion of each cycle it will not be far ahead of Neptune coming behind it, and Neptune will start to decrease Pluto's period again. The whole cycle takes about 20,000 years to complete.
Other factors
Numerical studies have shown that over millions of years, the general nature of the alignment between the orbits of Pluto and Neptune does not change. There are several other resonances and interactions that enhance Pluto's stability. These arise principally from two additional mechanisms (besides the 2:3 mean-motion resonance).
First, Pluto's argument of perihelion, the angle between the point where it crosses the ecliptic (or the invariant plane) and the point where it is closest to the Sun, librates around 90°. This means that when Pluto is closest to the Sun, it is at its farthest north of the plane of the Solar System, preventing encounters with Neptune. This is a consequence of the Kozai mechanism, which relates the eccentricity of an orbit to its inclination to a larger perturbing body—in this case, Neptune. Relative to Neptune, the amplitude of libration is 38°, and so the angular separation of Pluto's perihelion to the orbit of Neptune is always greater than 52° . The closest such angular separation occurs every 10,000 years.
Second, the longitudes of ascending nodes of the two bodies—the points where they cross the invariant plane—are in near-resonance with the above libration. When the two longitudes are the same—that is, when one could draw a straight line through both nodes and the Sun—Pluto's perihelion lies exactly at 90°, and hence it comes closest to the Sun when it is furthest north of Neptune's orbit. This is known as the 1:1 superresonance. All the Jovian planets (Jupiter, Saturn, Uranus, and Neptune) play a role in the creation of the superresonance.
Orcus
The second-largest known plutino, Orcus, has a diameter around 900 km and is in a very similar orbit to that of Pluto. However, the orbits of Pluto and Orcus are out of phase, so that the two never approach each other. It has been termed the "anti-Pluto", and is named for the Etruscan counterpart to the god Pluto.
Rotation
Pluto's rotation period, its day, is equal to 6.387 Earth days. Like Uranus and 2 Pallas, Pluto rotates on its "side" in its orbital plane, with an axial tilt of 120°, and so its seasonal variation is extreme; at its solstices, one-fourth of its surface is in continuous daylight, whereas another fourth is in continuous darkness. The reason for this unusual orientation has been debated. Research from the University of Arizona has suggested that it may be due to the way that a body's spin will always adjust to minimize energy. This could mean a body reorienting itself to put extraneous mass near the equator and regions lacking mass tend towards the poles. This is called polar wander. According to a paper released from the University of Arizona, this could be caused by masses of frozen nitrogen building up in shadowed areas of the dwarf planet. These masses would cause the body to reorient itself, leading to its unusual axial tilt of 120°. The buildup of nitrogen is due to Pluto's vast distance from the Sun. At the equator, temperatures can drop to , causing nitrogen to freeze as water would freeze on Earth. The same polar wandering effect seen on Pluto would be observed on Earth were the Antarctic ice sheet several times larger.
Geology
Surface
The plains on Pluto's surface are composed of more than 98 percent nitrogen ice, with traces of methane and carbon monoxide. Nitrogen and carbon monoxide are most abundant on the anti-Charon face of Pluto (around 180° longitude, where Tombaugh Regio's western lobe, Sputnik Planitia, is located), whereas methane is most abundant near 300° east. The mountains are made of water ice. Pluto's surface is quite varied, with large differences in both brightness and color. Pluto is one of the most contrastive bodies in the Solar System, with as much contrast as Saturn's moon Iapetus. The color varies from charcoal black, to dark orange and white. Pluto's color is more similar to that of Io with slightly more orange and significantly less red than Mars. Notable geographical features include Tombaugh Regio, or the "Heart" (a large bright area on the side opposite Charon), Belton Regio, or the "Whale" (a large dark area on the trailing hemisphere), and the "Brass Knuckles" (a series of equatorial dark areas on the leading hemisphere).
Sputnik Planitia, the western lobe of the "Heart", is a 1,000 km-wide basin of frozen nitrogen and carbon monoxide ices, divided into polygonal cells, which are interpreted as convection cells that carry floating blocks of water ice crust and sublimation pits towards their margins; there are obvious signs of glacial flows both into and out of the basin. It has no craters that were visible to New Horizons, indicating that its surface is less than 10 million years old. Latest studies have shown that the surface has an age of years.
The New Horizons science team summarized initial findings as "Pluto displays a surprisingly wide variety of geological landforms, including those resulting from glaciological and surface–atmosphere interactions as well as impact, tectonic, possible cryovolcanic, and mass-wasting processes."
In Western parts of Sputnik Planitia there are fields of transverse dunes formed by the winds blowing from the center of Sputnik Planitia in the direction of surrounding mountains. The dune wavelengths are in the range of 0.4–1 km and likely consist of methane particles 200–300 μm in size.
Internal structure
Pluto's density is . Because the decay of radioactive elements would eventually heat the ices enough for the rock to separate from them, scientists expect that Pluto's internal structure is differentiated, with the rocky material having settled into a dense core surrounded by a mantle of water ice. The pre–New Horizons estimate for the diameter of the core is , 70% of Pluto's diameter.
It is possible that such heating continues, creating a subsurface ocean of liquid water thick at the core–mantle boundary. In September 2016, scientists at Brown University simulated the impact thought to have formed Sputnik Planitia, and showed that it might have been the result of liquid water upweling from below after the collision, implying the existence of a subsurface ocean at least 100 km deep. In June 2020, astronomers reported evidence that Pluto may have had a subsurface ocean, and consequently may have been habitable, when it was first formed. In March 2022, a team of researchers proposed that the mountains Wright Mons and Piccard Mons are actually a merger of many smaller cryovolcanic domes, suggesting a source of heat on the body at levels previously thought not possible.
Mass and size
Pluto's diameter is and its mass is , 17.7% that of the Moon (0.22% that of Earth). Its surface area is , or just slightly bigger than Russia or Antarctica (particularly including the Antarctic sea ice during winter). Its surface gravity is 0.063 g (compared to 1 g for Earth and 0.17 g for the Moon). This gives Pluto an escape velocity of 4,363.2 km per hour / 2,711.167 miles per hour (as compared to Earth's 40,270 km per hour / 25,020 miles per hour). Pluto is more than twice the diameter and a dozen times the mass of Ceres, the largest object in the asteroid belt. It is less massive than the dwarf planet Eris, a trans-Neptunian object discovered in 2005, though Pluto has a larger diameter of 2,376.6 km compared to Eris's approximate diameter of 2,326 km.
With less than 0.2 lunar masses, Pluto is much less massive than the terrestrial planets, and also less massive than seven moons: Ganymede, Titan, Callisto, Io, the Moon, Europa, and Triton. The mass is much less than thought before Charon was discovered.
The discovery of Pluto's satellite Charon in 1978 enabled a determination of the mass of the Pluto–Charon system by application of Newton's formulation of Kepler's third law. Observations of Pluto in occultation with Charon allowed scientists to establish Pluto's diameter more accurately, whereas the invention of adaptive optics allowed them to determine its shape more accurately.
Determinations of Pluto's size have been complicated by its atmosphere and hydrocarbon haze. In March 2014, Lellouch, de Bergh et al. published findings regarding methane mixing ratios in Pluto's atmosphere consistent with a Plutonian diameter greater than 2,360 km, with a "best guess" of 2,368 km. On July 13, 2015, images from NASA's New Horizons mission Long Range Reconnaissance Imager (LORRI), along with data from the other instruments, determined Pluto's diameter to be , which was later revised to be on July 24, and later to . Using radio occultation data from the New Horizons Radio Science Experiment (REX), the diameter was found to be .
Atmosphere
Pluto has a tenuous atmosphere consisting of nitrogen (N2), methane (CH4), and carbon monoxide (CO), which are in equilibrium with their ices on Pluto's surface. According to the measurements by New Horizons, the surface pressure is about 1 Pa (10 μbar), roughly one million to 100,000 times less than Earth's atmospheric pressure. It was initially thought that, as Pluto moves away from the Sun, its atmosphere should gradually freeze onto the surface; studies of New Horizons data and ground-based occultations show that Pluto's atmospheric density increases, and that it likely remains gaseous throughout Pluto's orbit. New Horizons observations showed that atmospheric escape of nitrogen to be 10,000 times less than expected. Alan Stern has contended that even a small increase in Pluto's surface temperature can lead to exponential increases in Pluto's atmospheric density; from 18 hPa to as much as 280 hPa (three times that of Mars to a quarter that of the Earth). At such densities, nitrogen could flow across the surface as liquid. Just like sweat cools the body as it evaporates from the skin, the sublimation of Pluto's atmosphere cools its surface. Pluto has no or almost no troposphere; observations by New Horizons suggest only a thin tropospheric boundary layer. Its thickness in the place of measurement was 4 km, and the temperature was 37±3 K. The layer is not continuous.
In July 2019, an occultation by Pluto showed that its atmospheric pressure, against expectations, had fallen by 20% since 2016. In 2021, astronomers at the Southwest Research Institute confirmed the result using data from an occultation in 2018, which showed that light was appearing less gradually from behind Pluto's disc, indicating a thinning atmosphere.
The presence of methane, a powerful greenhouse gas, in Pluto's atmosphere creates a temperature inversion, with the average temperature of its atmosphere tens of degrees warmer than its surface, though observations by New Horizons have revealed Pluto's upper atmosphere to be far colder than expected (70 K, as opposed to about 100 K). Pluto's atmosphere is divided into roughly 20 regularly spaced haze layers up to 150 km high, thought to be the result of pressure waves created by airflow across Pluto's mountains.
Natural satellites
Pluto has five known natural satellites. The largest and closest to Pluto is Charon. First identified in 1978 by astronomer James Christy, Charon is the only moon of Pluto that may be in hydrostatic equilibrium. Charon's mass is sufficient to cause the barycenter of the Pluto–Charon system to be outside Pluto. Beyond Charon there are four much smaller circumbinary moons. In order of distance from Pluto they are Styx, Nix, Kerberos, and Hydra. Nix and Hydra were both discovered in 2005, Kerberos was discovered in 2011, and Styx was discovered in 2012. The satellites' orbits are circular (eccentricity < 0.006) and coplanar with Pluto's equator (inclination < 1°), and therefore tilted approximately 120° relative to Pluto's orbit. The Plutonian system is highly compact: the five known satellites orbit within the inner 3% of the region where prograde orbits would be stable.
The orbital periods of all Pluto's moons are linked in a system of orbital resonances and near-resonances. When precession is accounted for, the orbital periods of Styx, Nix, and Hydra are in an exact 18:22:33 ratio. There is a sequence of approximate ratios, 3:4:5:6, between the periods of Styx, Nix, Kerberos, and Hydra with that of Charon; the ratios become closer to being exact the further out the moons are.
The Pluto–Charon system is one of the few in the Solar System whose barycenter lies outside the primary body; the Patroclus–Menoetius system is a smaller example, and the Sun–Jupiter system is the only larger one. The similarity in size of Charon and Pluto has prompted some astronomers to call it a double dwarf planet. The system is also unusual among planetary systems in that each is tidally locked to the other, which means that Pluto and Charon always have the same hemisphere facing each other — a property shared by only one other known system, Eris and Dysnomia. From any position on either body, the other is always at the same position in the sky, or always obscured. This also means that the rotation period of each is equal to the time it takes the entire system to rotate around its barycenter.
Pluto's moons are hypothesized to have been formed by a collision between Pluto and a similar-sized body, early in the history of the Solar System. The collision released material that consolidated into the moons around Pluto.
Quasi-satellite
In 2012, it was calculated that 15810 Arawn could be a quasi-satellite of Pluto, a specific type of co-orbital configuration. According to the calculations, the object would be a quasi-satellite of Pluto for about 350,000 years out of every two-million-year period. Measurements made by the New Horizons spacecraft in 2015 made it possible to calculate the orbit of Arawn more accurately, and confirmed the earlier ones. However, it is not agreed upon among astronomers whether Arawn should be classified as a quasi-satellite of Pluto based on its orbital dynamics, since its orbit is primarily controlled by Neptune with only occasional perturbations by Pluto.
Origin
Pluto's origin and identity had long puzzled astronomers. One early hypothesis was that Pluto was an escaped moon of Neptune knocked out of orbit by Neptune's largest moon, Triton. This idea was eventually rejected after dynamical studies showed it to be impossible because Pluto never approaches Neptune in its orbit.
Pluto's true place in the Solar System began to reveal itself only in 1992, when astronomers began to find small icy objects beyond Neptune that were similar to Pluto not only in orbit but also in size and composition. This trans-Neptunian population is thought to be the source of many short-period comets. Pluto is the largest member of the Kuiper belt, a stable belt of objects located between 30 and 50 AU from the Sun. As of 2011, surveys of the Kuiper belt to magnitude 21 were nearly complete and any remaining Pluto-sized objects are expected to be beyond 100 AU from the Sun. Like other Kuiper-belt objects (KBOs), Pluto shares features with comets; for example, the solar wind is gradually blowing Pluto's surface into space. It has been claimed that if Pluto were placed as near to the Sun as Earth, it would develop a tail, as comets do. This claim has been disputed with the argument that Pluto's escape velocity is too high for this to happen. It has been proposed that Pluto may have formed as a result of the agglomeration of numerous comets and Kuiper-belt objects.
Though Pluto is the largest Kuiper belt object discovered, Neptune's moon Triton, which is larger than Pluto, is similar to it both geologically and atmospherically, and is thought to be a captured Kuiper belt object. Eris (see above) is about the same size as Pluto (though more massive) but is not strictly considered a member of the Kuiper belt population. Rather, it is considered a member of a linked population called the scattered disc.
Like other members of the Kuiper belt, Pluto is thought to be a residual planetesimal; a component of the original protoplanetary disc around the Sun that failed to fully coalesce into a full-fledged planet. Most astronomers agree that Pluto owes its position to a sudden migration undergone by Neptune early in the Solar System's formation. As Neptune migrated outward, it approached the objects in the proto-Kuiper belt, setting one in orbit around itself (Triton), locking others into resonances, and knocking others into chaotic orbits. The objects in the scattered disc, a dynamically unstable region overlapping the Kuiper belt, are thought to have been placed in their positions by interactions with Neptune's migrating resonances. A computer model created in 2004 by Alessandro Morbidelli of the Observatoire de la Côte d'Azur in Nice suggested that the migration of Neptune into the Kuiper belt may have been triggered by the formation of a 1:2 resonance between Jupiter and Saturn, which created a gravitational push that propelled both Uranus and Neptune into higher orbits and caused them to switch places, ultimately doubling Neptune's distance from the Sun. The resultant expulsion of objects from the proto-Kuiper belt could also explain the Late Heavy Bombardment 600 million years after the Solar System's formation and the origin of the Jupiter trojans. It is possible that Pluto had a near-circular orbit about 33 AU from the Sun before Neptune's migration perturbed it into a resonant capture. The Nice model requires that there were about a thousand Pluto-sized bodies in the original planetesimal disk, which included Triton and Eris.
Observation and exploration
Observation
Pluto's distance from Earth makes its in-depth study and exploration difficult. Pluto's visual apparent magnitude averages 15.1, brightening to 13.65 at perihelion. To see it, a telescope is required; around 30 cm (12 in) aperture being desirable. It looks star-like and without a visible disk even in large telescopes, because its angular diameter is maximum 0.11".
The earliest maps of Pluto, made in the late 1980s, were brightness maps created from close observations of eclipses by its largest moon, Charon. Observations were made of the change in the total average brightness of the Pluto–Charon system during the eclipses. For example, eclipsing a bright spot on Pluto makes a bigger total brightness change than eclipsing a dark spot. Computer processing of many such observations can be used to create a brightness map. This method can also track changes in brightness over time.
Better maps were produced from images taken by the Hubble Space Telescope (HST), which offered higher resolution, and showed considerably more detail, resolving variations several hundred kilometers across, including polar regions and large bright spots. These maps were produced by complex computer processing, which finds the best-fit projected maps for the few pixels of the Hubble images. These remained the most detailed maps of Pluto until the flyby of New Horizons in July 2015, because the two cameras on the HST used for these maps were no longer in service.
Exploration
The New Horizons spacecraft, which flew by Pluto in July 2015, is the first and so far only attempt to explore Pluto directly. Launched in 2006, it captured its first (distant) images of Pluto in late September 2006 during a test of the Long Range Reconnaissance Imager. The images, taken from a distance of approximately 4.2 billion kilometers, confirmed the spacecraft's ability to track distant targets, critical for maneuvering toward Pluto and other Kuiper belt objects. In early 2007 the craft made use of a gravity assist from Jupiter.New Horizons made its closest approach to Pluto on July 14, 2015, after a 3,462-day journey across the Solar System. Scientific observations of Pluto began five months before the closest approach and continued for at least a month after the encounter. Observations were conducted using a remote sensing package that included imaging instruments and a radio science investigation tool, as well as spectroscopic and other experiments. The scientific goals of New Horizons were to characterize the global geology and morphology of Pluto and its moon Charon, map their surface composition, and analyze Pluto's neutral atmosphere and its escape rate. On October 25, 2016, at 05:48 pm ET, the last bit of data (of a total of 50 billion bits of data; or 6.25 gigabytes) was received from New Horizons from its close encounter with Pluto.
Since the New Horizons flyby, scientists have advocated for an orbiter mission that would return to Pluto to fulfill new science objectives. They include mapping the surface at per pixel, observations of Pluto's smaller satellites, observations of how Pluto changes as it rotates on its axis, investigations of a possible subsurface ocean, and topographic mapping of Pluto's regions that are covered in long-term darkness due to its axial tilt. The last objective could be accomplished using laser pulses to generate a complete topographic map of Pluto. New Horizons principal investigator Alan Stern has advocated for a Cassini-style orbiter that would launch around 2030 (the 100th anniversary of Pluto's discovery) and use Charon's gravity to adjust its orbit as needed to fulfill science objectives after arriving at the Pluto system. The orbiter could then use Charon's gravity to leave the Pluto system and study more KBOs after all Pluto science objectives are completed. A conceptual study funded by the NASA Innovative Advanced Concepts (NIAC) program describes a fusion-enabled Pluto orbiter and lander based on the Princeton field-reversed configuration reactor.Fusion-Enabled Pluto Orbiter and Lander – Phase I Final Report . (PDF) Stephanie Thomas, Princeton Satellite Systems. 2017.New Horizons imaged all of Pluto's northern hemisphere, and the equatorial regions down to about 30° South. Higher southern latitudes have only been observed, at very low resolution, from Earth. Images from the Hubble Space Telescope in 1996 cover 85% of Pluto and show large albedo features down to about 75° South. This is enough to show the extent of the temperate-zone maculae. Later images had slightly better resolution, due to minor improvements in Hubble instrumentation. The equatorial region of the sub-Charon hemisphere of Pluto has only been imaged at low resolution, as New Horizons made its closest approach to the anti-Charon hemisphere.
Some albedo variations in the higher southern latitudes could be detected by New Horizons using Charon-shine (light reflected off Charon). The south polar region seems to be darker than the north polar region, but there is a high-albedo region in the southern hemisphere that may be a regional nitrogen or methane ice deposit.
| Physical sciences | Solar System | null |
44474 | https://en.wikipedia.org/wiki/Saturn | Saturn | Saturn is the sixth planet from the Sun and the second largest in the Solar System, after Jupiter. It is a gas giant, with an average radius of about nine times that of Earth. It has an eighth the average density of Earth, but is over 95 times more massive. Even though Saturn is almost as big as Jupiter, Saturn has less than a third its mass. Saturn orbits the Sun at a distance of , with an orbital period of 29.45 years.
Saturn's interior is thought to be composed of a rocky core, surrounded by a deep layer of metallic hydrogen, an intermediate layer of liquid hydrogen and liquid helium, and an outer layer of gas. Saturn has a pale yellow hue, due to ammonia crystals in its upper atmosphere. An electrical current in the metallic hydrogen layer is thought to give rise to Saturn's planetary magnetic field, which is weaker than Earth's, but has a magnetic moment 580 times that of Earth because of Saturn's greater size. Saturn's magnetic field strength is about a twentieth that of Jupiter. The outer atmosphere is generally bland and lacking in contrast, although long-lived features can appear. Wind speeds on Saturn can reach .
The planet has a bright and extensive system of rings, composed mainly of ice particles, with a smaller amount of rocky debris and dust. At least 146 moons orbit the planet, of which 63 are officially named; these do not include the hundreds of moonlets in the rings. Titan, Saturn's largest moon and the second largest in the Solar System, is larger (and less massive) than the planet Mercury and is the only moon in the Solar System that has a substantial atmosphere.
Name and symbol
Saturn is named after the Roman god of wealth and agriculture, who was the father of the god Jupiter. Its astronomical symbol has been traced back to the Greek Oxyrhynchus Papyri, where it can be seen to be a Greek kappa-rho ligature with a horizontal stroke, as an abbreviation for Κρονος (Cronus), the Greek name for the planet (). It later came to look like a lower-case Greek eta, with the cross added at the top in the 16th century to Christianize this pagan symbol.
The Romans named the seventh day of the week Saturday, Sāturni diēs, "Saturn's Day", for the planet Saturn.
Physical characteristics
Saturn is a gas giant, composed predominantly of hydrogen and helium. It lacks a definite surface, though it is likely to have a solid core. The planet's rotation makes it an oblate spheroid—a ball flattened at the poles and bulging at the equator. Its equatorial radius is more than 10% larger than the polar radius: 60,268 km versus 54,364 km (37,449 mi versus 33,780 mi). Jupiter, Uranus, and Neptune, the other giant planets in the Solar System, are less oblate. The combination of the bulge and the rotation rate means that the effective surface gravity along the equator, , is 74% of what it is at the poles and is lower than the surface gravity of Earth. However, the equatorial escape velocity, nearly , is much higher than that of Earth.
Saturn is the only planet of the Solar System that is less dense than water—about 30% less. Although Saturn's core is considerably denser than water, the average specific density of the planet is , because of the atmosphere. Jupiter has 318 times Earth's mass, and Saturn is 95 times Earth's mass. Together, Jupiter and Saturn hold 92% of the total planetary mass in the Solar System.
Internal structure
Despite consisting mostly of hydrogen and helium, most of Saturn's mass is not in the gas phase, because hydrogen becomes a non-ideal liquid when the density is above , which is reached at a radius containing 99.9% of Saturn's mass. The temperature, pressure, and density inside Saturn all rise steadily toward the core, which causes hydrogen to be a metal in the deeper layers.
Standard planetary models suggest that the interior of Saturn is similar to that of Jupiter, having a small rocky core surrounded by hydrogen and helium, with trace amounts of various volatiles. Analysis of the distortion shows that Saturn is substantially more centrally condensed than Jupiter and therefore contains much more material denser than hydrogen near its center. Saturn's central regions are about 50% hydrogen by mass, and Jupiter's are about 67% hydrogen.
This core is similar in composition to Earth, but is more dense. The examination of Saturn's gravitational moment, in combination with physical models of the interior, has allowed constraints to be placed on the mass of Saturn's core. In 2004, scientists estimated that the core must be 9–22 times the mass of Earth, which corresponds to a diameter of about . However, measurements of Saturn's rings suggest a much more diffuse core, with a mass equal to about 17 Earths and a radius equal to about 60% of Saturn's entire radius. This is surrounded by a thicker, liquid metallic hydrogen layer, followed by a liquid layer of helium-saturated molecular hydrogen, which gradually transitions to a gas as altitude increases. The outermost layer spans about and consists of gas.
Saturn has a hot interior, reaching at its core, and radiates 2.5 times more energy into space than it receives from the Sun. Jupiter's thermal energy is generated by the Kelvin–Helmholtz mechanism of slow gravitational compression; but such a process alone may not be sufficient to explain heat production for Saturn, because it is less massive. An alternative or additional mechanism may be the generation of heat through the "raining out" of droplets of helium deep in Saturn's interior. As the droplets descend through the lower-density hydrogen, the process releases heat by friction and leaves Saturn's outer layers depleted of helium. These descending droplets may have accumulated into a helium shell surrounding the core. Rainfalls of diamonds have been suggested to occur within Saturn, as well as in Jupiter and ice giants Uranus and Neptune.
Atmosphere
The outer atmosphere of Saturn contains 96.3% molecular hydrogen and 3.25% helium by volume. The proportion of helium is significantly deficient compared to the abundance of this element in the Sun. The quantity of elements heavier than helium (metallicity) is not known precisely, but the proportions are assumed to match the primordial abundances from the formation of the Solar System. The total mass of these heavier elements is estimated to be 19–31 times the mass of Earth, with a significant fraction located in Saturn's core region.
Trace amounts of ammonia, acetylene, ethane, propane, phosphine, and methane have been detected in Saturn's atmosphere. The upper clouds are composed of ammonia crystals, while the lower level clouds appear to consist of either ammonium hydrosulfide () or water. Ultraviolet radiation from the Sun causes methane photolysis in the upper atmosphere, leading to a series of hydrocarbon chemical reactions with the resulting products being carried downward by eddies and diffusion. This photochemical cycle is modulated by Saturn's annual seasonal cycle. Cassini observed a series of cloud features found in northern latitudes, nicknamed the "String of Pearls". These features are cloud clearings that reside in deeper cloud layers.
Cloud layers
Saturn's atmosphere exhibits a banded pattern similar to Jupiter's, but Saturn's bands are much fainter and are much wider near the equator. The nomenclature used to describe these bands is the same as on Jupiter. Saturn's finer cloud patterns were not observed until the flybys of the Voyager spacecraft during the 1980s. Since then, Earth-based telescopy has improved to the point where regular observations can be made.
The composition of the clouds varies with depth and increasing pressure. In the upper cloud layers, with temperatures in the range of 100–160 K and pressures extending between 0.5–2 bar, the clouds consist of ammonia ice. Water ice clouds begin at a level where the pressure is about 2.5 bar and extend down to 9.5 bar, where temperatures range from 185 to 270 K. Intermixed in this layer is a band of ammonium hydrosulfide ice, lying in the pressure range 3–6 bar with temperatures of 190–235 K. Finally, the lower layers, where pressures are between 10 and 20 bar and temperatures are 270–330 K, contains a region of water droplets with ammonia in aqueous solution.
Saturn's usually bland atmosphere occasionally exhibits long-lived ovals and other features common on Jupiter. In 1990, the Hubble Space Telescope imaged an enormous white cloud near Saturn's equator that was not present during the Voyager encounters, and in 1994 another smaller storm was observed. The 1990 storm was an example of a Great White Spot, a short-lived phenomenon that occurs once every Saturnian year, roughly every 30 Earth years, around the time of the northern hemisphere's summer solstice. Previous Great White Spots were observed in 1876, 1903, 1933, and 1960, with the 1933 storm being the best observed. The latest giant storm was observed in 2010. In 2015, researchers used Very Large Array telescope to study Saturnian atmosphere, and reported that they found "long-lasting signatures of all mid-latitude giant storms, a mixture of equatorial storms up to hundreds of years old, and potentially an unreported older storm at 70°N".
The winds on Saturn are the second fastest among the Solar System's planets, after Neptune's. Voyager data indicate peak easterly winds of . In images from the Cassini spacecraft during 2007, Saturn's northern hemisphere displayed a bright blue hue, similar to Uranus. The color was most likely caused by Rayleigh scattering. Thermography has shown that Saturn's south pole has a warm polar vortex, the only known example of such a phenomenon in the Solar System. Whereas temperatures on Saturn are normally −185 °C, temperatures on the vortex often reach as high as −122 °C, suspected to be the warmest spot on Saturn.
Hexagonal cloud patterns
A persisting hexagonal wave pattern around the north polar vortex in the atmosphere at about 78°N was first noted in the Voyager images. The sides of the hexagon are each about long, which is longer than the diameter of the Earth. The entire structure rotates with a period of (the same period as that of the planet's radio emissions) which is assumed to be equal to the period of rotation of Saturn's interior. The hexagonal feature does not shift in longitude like the other clouds in the visible atmosphere. The pattern's origin is a matter of much speculation. Most scientists think it is a standing wave pattern in the atmosphere. Polygonal shapes have been replicated in the laboratory through differential rotation of fluids.
HST imaging of the south polar region indicates the presence of a jet stream, but no strong polar vortex nor any hexagonal standing wave. NASA reported in November 2006 that Cassini had observed a "hurricane-like" storm locked to the south pole that had a clearly defined eyewall. Eyewall clouds had not previously been seen on any planet other than Earth. For example, images from the Galileo spacecraft did not show an eyewall in the Great Red Spot of Jupiter.
The south pole storm may have been present for billions of years. This vortex is comparable to the size of Earth, and it has winds of 550 km/h.
Magnetosphere
Saturn has an intrinsic magnetic field that has a simple, symmetric shape—a magnetic dipole. Its strength at the equator—0.2 gauss (20 μT)—is approximately one twentieth of that of the field around Jupiter and slightly weaker than Earth's magnetic field. As a result, Saturn's magnetosphere is much smaller than Jupiter's.
When Voyager 2 entered the magnetosphere, the solar wind pressure was high and the magnetosphere extended only 19 Saturn radii, or 1.1 million km (684,000 mi), although it enlarged within several hours, and remained so for about three days. Most probably, the magnetic field is generated similarly to that of Jupiter—by currents in the liquid metallic-hydrogen layer called a metallic-hydrogen dynamo. This magnetosphere is efficient at deflecting the solar wind particles from the Sun. The moon Titan orbits within the outer part of Saturn's magnetosphere and contributes plasma from the ionized particles in Titan's outer atmosphere. Saturn's magnetosphere, like Earth's, produces aurorae.
Orbit and rotation
The average distance between Saturn and the Sun is over 1.4 billion kilometers (9 AU). With an average orbital speed of 9.68 km/s, it takes Saturn 10,759 Earth days (or about years) to finish one revolution around the Sun. As a consequence, it forms a near 5:2 mean-motion resonance with Jupiter. The elliptical orbit of Saturn is inclined 2.48° relative to the orbital plane of the Earth. The perihelion and aphelion distances are, respectively, 9.195 and 9.957 AU, on average. The visible features on Saturn rotate at different rates depending on latitude, and multiple rotation periods have been assigned to various regions (as in Jupiter's case).
Astronomers use three different systems for specifying the rotation rate of Saturn. System I has a period of (844.3°/d) and encompasses the Equatorial Zone, the South Equatorial Belt, and the North Equatorial Belt. The polar regions are considered to have rotation rates similar to System I. All other Saturnian latitudes, excluding the north and south polar regions, are indicated as System II and have been assigned a rotation period of (810.76°/d). System III refers to Saturn's internal rotation rate. Based on radio emissions from the planet detected by Voyager 1 and Voyager 2, System III has a rotation period of (810.8°/d). System III has largely superseded System II.
A precise value for the rotation period of the interior remains elusive. While approaching Saturn in 2004, Cassini found that the radio rotation period of Saturn had increased appreciably, to approximately . An estimate of Saturn's rotation (as an indicated rotation rate for Saturn as a whole) based on a compilation of various measurements from the Cassini, Voyager, and Pioneer probes is . Studies of the planet's C Ring yield a rotation period of .
In March 2007, it was found that the variation in radio emissions from the planet did not match Saturn's rotation rate. This variance may be caused by geyser activity on Saturn's moon Enceladus. The water vapor emitted into Saturn's orbit by this activity becomes charged and creates a drag upon Saturn's magnetic field, slowing its rotation slightly relative to the rotation of the planet.
An apparent oddity for Saturn is that it does not have any known trojan asteroids. These are minor planets that orbit the Sun at the stable Lagrangian points, designated L4 and L5, located at 60° angles to the planet along its orbit. Trojan asteroids have been discovered for Mars, Jupiter, Uranus, and Neptune. Orbital resonance mechanisms, including secular resonance, are believed to be the cause of the missing Saturnian trojans.
Natural satellites
Saturn has 146 known moons, 63 of which have formal names. It is estimated that there are another outer irregular moons larger than in diameter. In addition, there is evidence of dozens to hundreds of moonlets with diameters of 40–500 meters in Saturn's rings, which are not considered to be true moons. Titan, the largest moon, comprises more than 90% of the mass in orbit around Saturn, including the rings. Saturn's second-largest moon, Rhea, may have a tenuous ring system of its own, along with a tenuous atmosphere.
Many of the other moons are small: 131 are less than 50 km in diameter. Traditionally, most of Saturn's moons have been named after Titans of Greek mythology. Titan is the only satellite in the Solar System with a major atmosphere, in which a complex organic chemistry occurs. It is the only satellite with hydrocarbon lakes.
On 6 June 2013, scientists at the IAA-CSIC reported the detection of polycyclic aromatic hydrocarbons in the upper atmosphere of Titan, a possible precursor for life. On 23 June 2014, NASA claimed to have strong evidence that nitrogen in the atmosphere of Titan came from materials in the Oort cloud, associated with comets, and not from the materials that formed Saturn in earlier times.
Saturn's moon Enceladus, which seems similar in chemical makeup to comets, has often been regarded as a potential habitat for microbial life. Evidence of this possibility includes the satellite's salt-rich particles having an "ocean-like" composition that indicates most of Enceladus's expelled ice comes from the evaporation of liquid salt water. A 2015 flyby by Cassini through a plume on Enceladus found most of the ingredients to sustain life forms that live by methanogenesis.
In April 2014, NASA scientists reported the possible beginning of a new moon within the A Ring, which was imaged by Cassini on 15 April 2013.
Planetary rings
Saturn is probably best known for the system of planetary rings that makes it visually unique. The rings extend from outward from Saturn's equator and average approximately in thickness. They are composed predominantly of water ice, with trace amounts of tholin impurities and a peppered coating of approximately 7% amorphous carbon. The particles that make up the rings range in size from specks of dust up to 10 m. While the other gas giants also have ring systems, Saturn's is the largest and most visible.
There is a debate on the age of the rings. One side supports that they are ancient, and were created simultaneously with Saturn from the original nebular material (around 4.6 billion years ago), or shortly after the LHB (around 4.1 to 3.8 billion years ago). The other side supports that they are much younger, created around 100 million years ago. An MIT research team, supporting the latter theory, proposed that the rings are remnant of a destroyed moon of Saturn, named ″Chrysalis″.
Beyond the main rings, at a distance of 12 million km (7.5 million mi) from the planet is the sparse Phoebe ring. It is tilted at an angle of 27° to the other rings and, like Phoebe, orbits in retrograde fashion.
Some of the moons of Saturn, including Pandora and Prometheus, act as shepherd moons to confine the rings and prevent them from spreading out. Pan and Atlas cause weak, linear density waves in Saturn's rings that have yielded more reliable calculations of their masses.
History of observation and exploration
The observation and exploration of Saturn can be divided into three phases: (1) pre-modern observations with the naked eye, (2) telescopic observations from Earth beginning in the 17th century, and (3) visitation by space probes, in orbit or on flyby. In the 21st century, telescopic observations continue from Earth (including Earth-orbiting observatories like the Hubble Space Telescope) and, until its 2017 retirement, from the Cassini orbiter around Saturn.
Pre-telescopic observation
Saturn has been known since prehistoric times, and in early recorded history it was a major character in various mythologies. Babylonian astronomers systematically observed and recorded the movements of Saturn. In ancient Greek, the planet was known as Phainon, and in Roman times it was known as the "star of Saturn" or the "star of the Sun (i.e. Helios)". In ancient Roman mythology, the planet Phainon was sacred to this agricultural god, from which the planet takes its modern name. The Romans considered the god Saturnus the equivalent of the Greek god Cronus; in modern Greek, the planet retains the name Cronus—: Kronos.
The Greek scientist Ptolemy based his calculations of Saturn's orbit on observations he made while it was in opposition. In Hindu astrology, there are nine astrological objects, known as Navagrahas. Saturn is known as "Shani" and judges everyone based on the good and bad deeds performed in life. Ancient Chinese and Japanese culture designated the planet Saturn as the "earth star" (). This was based on Five Elements which were traditionally used to classify natural elements.
In Hebrew, Saturn is called Shabbathai. Its angel is Cassiel. Its intelligence or beneficial spirit is 'Agȋȇl (), and its darker spirit (demon) is Zȃzȇl (). Zazel has been described as a great angel, invoked in Solomonic magic, who is "effective in love conjurations". In Ottoman Turkish, Urdu, and Malay, the name of Zazel is 'Zuhal', derived from the Arabic language ().
Telescopic pre-spaceflight observations
Saturn's rings require at least a 15-mm-diameter telescope to resolve and thus were not known to exist until Christiaan Huygens saw them in 1655 and published his observations in 1659. Galileo, with his primitive telescope in 1610, incorrectly thought of Saturn's appearing not quite round as two moons on Saturn's sides. It was not until Huygens used greater telescopic magnification that this notion was refuted, and the rings were truly seen for the first time. Huygens also discovered Saturn's moon Titan; Giovanni Domenico Cassini later discovered four other moons: Iapetus, Rhea, Tethys, and Dione. In 1675, Cassini discovered the gap now known as the Cassini Division.
No further discoveries of significance were made until 1789 when William Herschel discovered two further moons, Mimas and Enceladus. The irregularly shaped satellite Hyperion, which has a resonance with Titan, was discovered in 1848 by a British team.
In 1899, William Henry Pickering discovered Phoebe, a highly irregular satellite that does not rotate synchronously with Saturn as the larger moons do. Phoebe was the first such satellite found and it took more than a year to orbit Saturn in a retrograde orbit. During the early 20th century, research on Titan led to the confirmation in 1944 that it had a thick atmosphere—a feature unique among the Solar System's moons.
Spaceflight missions
Pioneer 11 flyby
Pioneer 11 made the first flyby of Saturn in September 1979, when it passed within of the planet's cloud tops. Images were taken of the planet and a few of its moons, although their resolution was too low to discern surface detail. The spacecraft also studied Saturn's rings, revealing the thin F-ring and the fact that dark gaps in the rings are bright when viewed at a high phase angle (towards the Sun), meaning that they contain fine light-scattering material. In addition, Pioneer 11 measured the temperature of Titan.
Voyager flybys
In November 1980, the Voyager 1 probe visited the Saturn system. It sent back the first high-resolution images of the planet, its rings and satellites. Surface features of various moons were seen for the first time. Voyager 1 performed a close flyby of Titan, increasing knowledge of the atmosphere of the moon. It proved that Titan's atmosphere is impenetrable at visible wavelengths; therefore no surface details were seen. The flyby changed the spacecraft's trajectory out of the plane of the Solar System.
Almost a year later, in August 1981, Voyager 2 continued the study of the Saturn system. More close-up images of Saturn's moons were acquired, as well as evidence of changes in the atmosphere and the rings. During the flyby, the probe's turnable camera platform stuck for a couple of days and some planned imaging was lost. Saturn's gravity was used to direct the spacecraft's trajectory towards Uranus.
The probes discovered and confirmed several new satellites orbiting near or within the planet's rings, as well as the small Maxwell Gap (a gap within the C Ring) and Keeler gap (a 42 km-wide gap in the A Ring).
Cassini–Huygens spacecraft
The Cassini–Huygens space probe entered orbit around Saturn on 1 July 2004. In June 2004, it conducted a close flyby of Phoebe, sending back high-resolution images and data. Cassini flyby of Saturn's largest moon, Titan, captured radar images of large lakes and their coastlines with numerous islands and mountains. The orbiter completed two Titan flybys before releasing the Huygens probe on 25 December 2004. Huygens descended onto the surface of Titan on 14 January 2005.
Starting in early 2005, scientists used Cassini to track lightning on Saturn. The power of the lightning is approximately 1,000 times that of lightning on Earth.
In 2006, NASA reported that Cassini had found evidence of liquid water reservoirs no more than tens of meters below the surface that erupt in geysers on Saturn's moon Enceladus. These jets of icy particles are emitted into orbit around Saturn from vents in the moon's south polar region. Over 100 geysers have been identified on Enceladus. In May 2011, NASA scientists reported that Enceladus "is emerging as the most habitable spot beyond Earth in the Solar System for life as we know it".
Cassini photographs have revealed a previously undiscovered planetary ring, outside the brighter main rings of Saturn and inside the G and E rings. The source of this ring is hypothesized to be the crashing of a meteoroid off Janus and Epimetheus. In July 2006, images were returned of hydrocarbon lakes near Titan's north pole, the presence of which were confirmed in January 2007. In March 2007, hydrocarbon seas were found near the North pole, the largest of which is almost the size of the Caspian Sea. In October 2006, the probe detected an diameter cyclone-like storm with an eyewall at Saturn's south pole.
From 2004 to 2 November 2009, the probe discovered and confirmed eight new satellites. In April 2013, Cassini sent back images of a hurricane at the planet's north pole 20 times larger than those found on Earth, with winds faster than . On 15 September 2017, the Cassini–Huygens spacecraft performed the "Grand Finale" of its mission: a number of passes through gaps between Saturn and Saturn's inner rings. The atmospheric entry of Cassini ended the mission.
Possible future missions
The continued exploration of Saturn is still considered to be a viable option for NASA as part of their ongoing New Frontiers program of missions. NASA previously requested for plans to be put forward for a mission to Saturn that included the Saturn Atmospheric Entry Probe, and possible investigations into the habitability and possible discovery of life on Saturn's moons Titan and Enceladus by Dragonfly.
Observation
Saturn is the most distant of the five planets easily visible to the naked eye from Earth, the other four being Mercury, Venus, Mars, and Jupiter. (Uranus, and occasionally 4 Vesta, are visible to the naked eye in dark skies.) Saturn appears to the naked eye in the night sky as a bright, yellowish point of light. The mean apparent magnitude of Saturn is 0.46 with a standard deviation of 0.34. Most of the magnitude variation is due to the inclination of the ring system relative to the Sun and Earth. The brightest magnitude, −0.55, occurs near the time when the plane of the rings is inclined most highly, and the faintest magnitude, 1.17, occurs around the time when they are least inclined. It takes approximately 29.4 years for the planet to complete an entire circuit of the ecliptic against the background constellations of the zodiac. Most people will require an optical aid (very large binoculars or a small telescope) that magnifies at least 30 times to achieve an image of Saturn's rings in which a clear resolution is present. When Earth passes through the ring plane, which occurs twice every Saturnian year (roughly every 15 Earth years), the rings briefly disappear from view because they are so thin. Such a "disappearance" will next occur in 2025, but Saturn will be too close to the Sun for observations.
Saturn and its rings are best seen when the planet is at, or near, opposition, the configuration of a planet when it is at an elongation of 180°, and thus appears opposite the Sun in the sky. A Saturnian opposition occurs every year—approximately every 378 days—and results in the planet appearing at its brightest. Both the Earth and Saturn orbit the Sun on eccentric orbits, which means their distances from the Sun vary over time, and therefore so do their distances from each other, hence varying the brightness of Saturn from one opposition to the next. Saturn also appears brighter when the rings are angled such that they are more visible. For example, during the opposition of 17 December 2002, Saturn appeared at its brightest due to the favorable orientation of its rings relative to the Earth, even though Saturn was closer to the Earth and Sun in late 2003.
From time to time, Saturn is occulted by the Moon (that is, the Moon covers up Saturn in the sky). As with all the planets in the Solar System, occultations of Saturn occur in "seasons". Saturnian occultations will take place monthly for about a 12-month period, followed by about a five-year period in which no such activity is registered. The Moon's orbit is inclined by several degrees relative to Saturn's, so occultations will only occur when Saturn is near one of the points in the sky where the two planes intersect (both the length of Saturn's year and the 18.6-Earth-year nodal precession period of the Moon's orbit influence the periodicity).
In science fiction
In Christopher Nolan's 2014 science fiction epic Interstellar, in proximity to Saturn is a wormhole leading to a planetary system in another galaxy, whose central object is a black hole known as Gargantua. The Endurance team enters the wormhole in the hopes of finding a habitable planet for humanity to settle as conditions on Earth deteriorate. At the end of the film, Cooper Station, named for the main character, is shown in orbit around Saturn.
| Physical sciences | Astronomy | null |
44475 | https://en.wikipedia.org/wiki/Uranus | Uranus | Uranus is the seventh planet from the Sun. It is a gaseous cyan-coloured ice giant. Most of the planet is made of water, ammonia, and methane in a supercritical phase of matter, which astronomy calls "ice" or volatiles. The planet's atmosphere has a complex layered cloud structure and has the lowest minimum temperature () of all the Solar System's planets. It has a marked axial tilt of 82.23° with a retrograde rotation period of 17 hours and 14 minutes. This means that in an 84-Earth-year orbital period around the Sun, its poles get around 42 years of continuous sunlight, followed by 42 years of continuous darkness.
Uranus has the third-largest diameter and fourth-largest mass among the Solar System's planets. Based on current models, inside its volatile mantle layer is a rocky core, and surrounding it is a thick hydrogen and helium atmosphere. Trace amounts of hydrocarbons (thought to be produced via hydrolysis) and carbon monoxide along with carbon dioxide (thought to have been originated from comets) have been detected in the upper atmosphere. There are many unexplained climate phenomena in Uranus's atmosphere, such as its peak wind speed of , variations in its polar cap, and its erratic cloud formation. The planet also has very low internal heat compared to other giant planets, the cause of which remains unclear.
Like the other giant planets, Uranus has a ring system, a magnetosphere, and many natural satellites. The extremely dark ring system reflects only about 2% of the incoming light. Uranus's 28 natural satellites include 18 known regular moons, of which 13 are small inner moons. Further out are the larger five major moons of the planet: Miranda, Ariel, Umbriel, Titania, and Oberon. Orbiting at a much greater distance from Uranus are the ten known irregular moons. The planet's magnetosphere is highly asymmetric and has many charged particles, which may be the cause of the darkening of its rings and moons.
Uranus is visible to the naked eye, but it is very dim and was not classified as a planet until 1781, when it was first observed by William Herschel. About seven decades after its discovery, consensus was reached that the planet be named after the Greek god Uranus (Ouranos), one of the Greek primordial deities. As of 2024, it had been visited up close only once when in 1986 the Voyager 2 probe flew by the planet. Though nowadays it can be resolved and observed by telescopes, there is much desire to revisit the planet, as shown by Planetary Science Decadal Survey's decision to make the proposed Uranus Orbiter and Probe mission a top priority in the 2023–2032 survey, and the CNSA's proposal to fly by the planet with a subprobe of Tianwen-4.
History
Like the classical planets, Uranus is visible to the naked eye, but it was never recognised as a planet by ancient observers because of its dimness and slow orbit. William Herschel first observed Uranus on 13 March 1781, leading to its discovery as a planet, expanding the known boundaries of the Solar System for the first time in history and making Uranus the first planet classified as such with the aid of a telescope. The discovery of Uranus also effectively doubled the size of the known Solar System because Uranus is around twice the distance from the Sun as the planet Saturn.
Discovery
Before its recognition as a planet, Uranus had been observed on numerous occasions, albeit generally misidentified as a star. The earliest possible known observation was by Hipparchus, who in 128 BC might have recorded it as a star for his star catalogue that was later incorporated into Ptolemy's Almagest. The earliest definite sighting was in 1690, when John Flamsteed observed it at least six times, cataloguing it as 34 Tauri. The French astronomer Pierre Charles Le Monnier observed Uranus at least twelve times between 1750 and 1769, including on four consecutive nights.
William Herschel observed Uranus on 13 March 1781 from the garden of his house at 19 New King Street in Bath, Somerset, England (now the Herschel Museum of Astronomy), and initially reported it (on 26 April 1781) as a comet. With a homemade 6.2-inch reflecting telescope, Herschel "engaged in a series of observations on the parallax of the fixed stars."
Herschel recorded in his journal: "In the quartile near ζ Tauri ... either [a] Nebulous star or perhaps a comet." On 17 March he noted: "I looked for the Comet or Nebulous Star and found that it is a Comet, for it has changed its place." When he presented his discovery to the Royal Society, he continued to assert that he had found a comet, but also implicitly compared it to a planet:
Herschel notified the Astronomer Royal Nevil Maskelyne of his discovery and received this flummoxed reply from him on 23 April 1781: "I don't know what to call it. It is as likely to be a regular planet moving in an orbit nearly circular to the sun as a Comet moving in a very eccentric ellipsis. I have not yet seen any coma or tail to it."
Although Herschel continued to describe his new object as a comet, other astronomers had already begun to suspect otherwise. Finnish-Swedish astronomer Anders Johan Lexell, working in Russia, was the first to compute the orbit of the new object. Its nearly circular orbit led him to the conclusion that it was a planet rather than a comet. Berlin astronomer Johann Elert Bode described Herschel's discovery as "a moving star that can be deemed a hitherto unknown planet-like object circulating beyond the orbit of Saturn". Bode concluded that its near-circular orbit was more like a planet's than a comet's.
The object was soon universally accepted as a new planet. By 1783, Herschel acknowledged this to Royal Society president Joseph Banks: "By the observation of the most eminent Astronomers in Europe it appears that the new star, which I had the honour of pointing out to them in March 1781, is a Primary Planet of our Solar System." In recognition of his achievement, King George III gave Herschel an annual stipend of £200 () on condition that he moved to Windsor so that the Royal Family could look through his telescopes.
Name
The name Uranus references the ancient Greek deity of the sky Uranus (), known as Caelus in Roman mythology, the father of Cronus (Saturn), grandfather of Zeus (Jupiter) and the great-grandfather of Ares (Mars), which was rendered as in Latin (). It is the only one of the eight planets whose English name derives from a figure of Greek mythology. The pronunciation of the name Uranus preferred among astronomers is , with the long "u" of English and stress on the first syllable as in Latin , in contrast to , with stress on the second syllable and a long a, though both are considered acceptable.
Consensus on the name was not reached until almost 70 years after the planet's discovery. During the original discussions following discovery, Maskelyne asked Herschel to "do the astronomical world the to give a name to your planet, which is entirely your own, [and] which we are so much obliged to you for the discovery of". In response to Maskelyne's request, Herschel decided to name the object (George's Star), or the "Georgian Planet" in honour of his new patron, King George III. He explained this decision in a letter to Joseph Banks:
Herschel's proposed name was not popular outside Britain and Hanover, and alternatives were soon proposed. Astronomer Jérôme Lalande proposed that it be named Herschel in honour of its discoverer. Swedish astronomer Erik Prosperin proposed the names Astraea, Cybele (now the names of asteroids), and Neptune, which would become the name of the next planet to be discovered. Georg Lichtenberg from Göttingen also supported Astraea (as Austräa), but she is traditionally associated with Virgo instead of Taurus. Neptune was supported by other astronomers who liked the idea of commemorating the victories of the British Royal Naval fleet in the course of the American Revolutionary War by calling the new planet either Neptune George III or Neptune Great Britain, a compromise Lexell suggested as well. Daniel Bernoulli suggested Hypercronius and Transaturnis. Minerva was also proposed.
In a March 1782 treatise, Johann Elert Bode proposed Uranus, the Latinised version of the Greek god of the sky, Ouranos. Bode argued that the name should follow the mythology so as not to stand out as different from the other planets, and that Uranus was an appropriate name as the father of the first generation of the Titans. He also noted the elegance of the name in that just as Saturn was the father of Jupiter, the new planet should be named after the father of Saturn. However, he was apparently unaware that Uranus was only the Latinised form of the deity's name, and the Roman equivalent was Caelus. In 1789, Bode's Royal Academy colleague Martin Klaproth named his newly discovered element uranium in support of Bode's choice. Ultimately, Bode's suggestion became the most widely used, and became universal in 1850 when HM Nautical Almanac Office, the final holdout, switched from using Georgium Sidus to Uranus.
Uranus has two astronomical symbols. The first to be proposed, , was proposed by Johann Gottfried Köhler at Bode's request in 1782. Köhler suggested that the new planet be given the symbol for platinum, which had been described scientifically only 30 years before. As there was no alchemical symbol for platinum, he suggested ⛢ or ⛢, a combination of the planetary-metal symbols ☉ (gold) and ♂ (iron), as platinum (or 'white gold') is found mixed with iron. Bode thought that an upright orientation, ⛢, fit better with the symbols for the other planets while remaining distinct. This symbol predominates in modern astronomical use in the rare cases that symbols are used at all. The second symbol, , was suggested by Lalande in 1784. In a letter to Herschel, Lalande described it as "" ("a globe surmounted by the first letter of your surname"). The second symbol is nearly universal in astrology.
In English-language popular culture, humour is often derived from the common pronunciation of Uranus's name, which resembles that of the phrase "your anus".
Uranus is called by a variety of names in other languages. Uranus's name is literally translated as the "Heavenly King star" in Chinese (), Japanese (天王星), Korean (천왕성), and Vietnamese (sao Thiên Vương). In Thai, its official name is (), as in English. Its other name in Thai is (, Star of Mṛtyu), after the Sanskrit word for 'death', (). In Mongolian, its name is (), translated as 'King of the Sky', reflecting its namesake god's role as the ruler of the heavens. In Hawaiian, its name is , the Hawaiian rendering of the name 'Herschel'.
Formation
It is argued that the differences between the ice giants and the gas giants arise from their formation history. The Solar System is hypothesised to have formed from a rotating disk of gas and dust known as the presolar nebula. Much of the nebula's gas, primarily hydrogen and helium, formed the Sun, and the dust grains collected together to form the first protoplanets. As the planets grew, some of them eventually accreted enough matter for their gravity to hold on to the nebula's leftover gas. The more gas they held onto, the larger they became; the larger they became, the more gas they held onto until a critical point was reached, and their size began to increase exponentially. The ice giants, with only a few Earth masses of nebular gas, never reached that critical point. Recent simulations of planetary migration have suggested that both ice giants formed closer to the Sun than their present positions, and moved outwards after formation (the Nice model).
Orbit and rotation
Uranus orbits the Sun once every 84 years. As viewed against the background of stars, since being discovered in 1781, the planet has returned to the point of its discovery northeast of the binary star Zeta Tauri twice—in March 1865 and March 1949—and will return to this location again in April 2033.
Its average distance from the Sun is roughly . The difference between its minimum and maximum distance from the Sun is 1.8 AU, larger than that of any other planet, though not as large as that of dwarf planet Pluto. The intensity of sunlight varies inversely with the square of the distance—on Uranus (at about 20 times the distance from the Sun compared to Earth), it is about 1/400 the intensity of light on Earth.
The orbital elements of Uranus were first calculated in 1783 by Pierre-Simon Laplace. With time, discrepancies began to appear between predicted and observed orbits, and in 1841, John Couch Adams first proposed that the differences might be due to the gravitational tug of an unseen planet. In 1845, Urbain Le Verrier began his own independent research into Uranus's orbit. On 23 September 1846, Johann Gottfried Galle located a new planet, later named Neptune, at nearly the position predicted by Le Verrier.
The rotational period of the interior of Uranus is 17 hours, 14 minutes. As on all giant planets, its upper atmosphere experiences strong winds in the direction of rotation. At some latitudes, such as about 60 degrees south, visible features of the atmosphere move much faster, making a full rotation in as little as 14 hours.
Axial tilt
The Uranian axis of rotation is approximately parallel to the plane of the Solar System, with an axial tilt of 82.23°. Depending on which pole is considered north, the tilt can be described either as 82.23° or as 97.8°. The former follows the International Astronomical Union definition that the north pole is the pole which lies on Earth's North's side of the invariable plane of the Solar System. Uranus has retrograde rotation when defined this way. Alternatively, the convention in which a body's north and south poles are defined according to the right-hand rule in relation to the direction of rotation, Uranus's axial tilt may be given instead as 97.8°, which reverses which pole is considered north and which is considered south and giving the planet prograde rotation. This gives it seasonal changes completely unlike those of the other planets. Pluto and asteroid 2 Pallas also have extreme axial tilts. Near the solstice, one pole faces the Sun continuously and the other faces away, with only a narrow strip around the equator experiencing a rapid day–night cycle, with the Sun low over the horizon. On the other side of Uranus's orbit, the orientation of the poles towards the Sun is reversed. Each pole gets around 42 years of continuous sunlight, followed by 42 years of darkness. Near the time of the equinoxes, the Sun faces the equator of Uranus, giving a period of day–night cycles similar to those seen on most of the other planets.
One result of this axis orientation is that, averaged over the Uranian year, the near-polar regions of Uranus receive a greater energy input from the Sun than its equatorial regions. Nevertheless, Uranus is hotter at its equator than at its poles. The underlying mechanism that causes this is unknown. The reason for Uranus's unusual axial tilt is also not known with certainty, but the usual speculation is that during the formation of the Solar System, an Earth-sized protoplanet collided with Uranus, causing the skewed orientation. Research by Jacob Kegerreis of Durham University suggests that the tilt resulted from a rock larger than Earth crashing into the planet 3 to 4 billion years ago. Uranus's south pole was pointed almost directly at the Sun at the time of Voyager 2 flyby in 1986.
Visibility from Earth
The mean apparent magnitude of Uranus is 5.68 with a standard deviation of 0.17, while the extremes are 5.38 and 6.03. This range of brightness is near the limit of naked eye visibility. Much of the variability is dependent upon the planetary latitudes being illuminated from the Sun and viewed from the Earth. Its angular diameter is between 3.4 and 3.7 arcseconds, compared with 16 to 20 arcseconds for Saturn and 32 to 45 arcseconds for Jupiter. At opposition, Uranus is visible to the naked eye in dark skies, and becomes an easy target even in urban conditions with binoculars. On larger amateur telescopes with an objective diameter of between 15 and 23 cm, Uranus appears as a pale cyan disk with distinct limb darkening. With a large telescope of 25 cm or wider, cloud patterns, as well as some of the larger satellites, such as Titania and Oberon, may be visible.
Internal structure
Uranus's mass is roughly 14.5 times that of Earth, making it the least massive of the giant planets. Its diameter is slightly larger than Neptune's at roughly four times that of Earth. A resulting density of 1.27 g/cm3 makes Uranus the second least dense planet, after Saturn. This value indicates that it is made primarily of various ices, such as water, ammonia, and methane. The total mass of ice in Uranus's interior is not precisely known, because different figures emerge depending on the model chosen; it must be between 9.3 and 13.5 Earth masses. Hydrogen and helium constitute only a small part of the total, with between 0.5 and 1.5 Earth masses. The remainder of the non-ice mass (0.5 to 3.7 Earth masses) is accounted for by rocky material.
The standard model of Uranus's structure is that it consists of three layers: a rocky (silicate/iron–nickel) core in the centre, an icy mantle in the middle, and an outer gaseous hydrogen/helium envelope. The core is relatively small, with a mass of only 0.55 Earth masses and a radius less than 20% of the planet; the mantle comprises its bulk, with around 13.4 Earth masses, and the upper atmosphere is relatively insubstantial, weighing about 0.5 Earth masses and extending for the last 20% of Uranus's radius. Uranus's core density is around 9 g/cm3, with a pressure in the centre of 8 million bars (800 GPa) and a temperature of about 5000 K. The ice mantle is not in fact composed of ice in the conventional sense, but of a hot and dense fluid consisting of water, ammonia and other volatiles. This fluid, which has a high electrical conductivity, is sometimes called a water–ammonia ocean.
The extreme pressure and temperature deep within Uranus may break up the methane molecules, with the carbon atoms condensing into crystals of diamond that rain down through the mantle like hailstones. This phenomenon is similar to diamond rains that are theorised by scientists to exist on Jupiter, Saturn, and Neptune. Very-high-pressure experiments at the Lawrence Livermore National Laboratory suggest that an ocean of metallic liquid carbon, perhaps with floating solid 'diamond-bergs', may comprise the base of the mantle.
The bulk compositions of Uranus and Neptune are different from those of Jupiter and Saturn, with ice dominating over gases, hence justifying their separate classification as ice giants. There may be a layer of ionic water where the water molecules break down into a soup of hydrogen and oxygen ions, and deeper down superionic water in which the oxygen crystallises but the hydrogen ions move freely within the oxygen lattice.
Although the model considered above is reasonably standard, it is not unique; other models also satisfy observations. For instance, if substantial amounts of hydrogen and rocky material are mixed in the ice mantle, the total mass of ices in the interior will be lower, and, correspondingly, the total mass of rocks and hydrogen will be higher. Presently available data does not allow a scientific determination of which model is correct. The fluid interior structure of Uranus means that it has no solid surface. The gaseous atmosphere gradually transitions into the internal liquid layers. For the sake of convenience, a revolving oblate spheroid set at the point at which atmospheric pressure equals 1 bar (100 kPa) is conditionally designated as a "surface". It has equatorial and polar radii of and , respectively. This surface is used throughout this article as a zero point for altitudes.
Internal heat
Uranus's internal heat appears markedly lower than that of the other giant planets; in astronomical terms, it has a low thermal flux. Why Uranus's internal temperature is so low is still not understood. Neptune, which is Uranus's near twin in size and composition, radiates 2.61 times as much energy into space as it receives from the Sun, but Uranus radiates hardly any excess heat at all. The total power radiated by Uranus in the far infrared (i.e. heat) part of the spectrum is times the solar energy absorbed in its atmosphere. Uranus's heat flux is only , which is lower than the internal heat flux of Earth of about . The lowest temperature recorded in Uranus's tropopause is , making Uranus the coldest planet in the Solar System.
One of the hypotheses for this discrepancy suggests the Earth-sized impactor theorised to be behind Uranus's axial tilt left the planet with a depleted core temperature, as the impact caused Uranus to expel most of its primordial heat. Another hypothesis is that some form of barrier exists in Uranus's upper layers that prevents the core's heat from reaching the surface. For example, convection may take place in a set of compositionally different layers, which may inhibit upward heat transport; perhaps double diffusive convection is a limiting factor.
In a 2021 study, the ice giants' interior conditions were mimicked by compressing water that contained minerals such as olivine and ferropericlase, thus showing that large amounts of magnesium could be dissolved in the liquid interiors of Uranus and Neptune. If Uranus has more of this magnesium than Neptune, it could form a thermal insulation layer, thus potentially explaining the planet's low temperature.
Atmosphere
Although there is no well-defined solid surface within Uranus's interior, the outermost part of Uranus's gaseous envelope that is accessible to remote sensing is called its atmosphere. Remote-sensing capability extends down to roughly 300 km below the level, with a corresponding pressure around and temperature of . The tenuous thermosphere extends over two planetary radii from the nominal surface, which is defined to lie at a pressure of 1 bar. The Uranian atmosphere can be divided into three layers: the troposphere, between altitudes of and pressures from 100 to 0.1 bar (10 MPa to 10 kPa); the stratosphere, spanning altitudes between and pressures of between (10 kPa to 10 μPa); and the thermosphere extending from 4,000 km to as high as 50,000 km from the surface. There is no mesosphere.
Composition
The composition of Uranus's atmosphere is different from its bulk, consisting mainly of molecular hydrogen and helium. The helium molar fraction, i.e. the number of helium atoms per molecule of gas, is in the upper troposphere, which corresponds to a mass fraction . This value is close to the protosolar helium mass fraction of , indicating that helium has not settled in its centre as it has in the gas giants. The third-most-abundant component of Uranus's atmosphere is methane (). Methane has prominent absorption bands in the visible and near-infrared (IR), making Uranus aquamarine or cyan in colour. Methane molecules account for 2.3% of the atmosphere by molar fraction below the methane cloud deck at the pressure level of ; this represents about 20 to 30 times the carbon abundance found in the Sun.
The mixing ratio is much lower in the upper atmosphere due to its extremely low temperature, which lowers the saturation level and causes excess methane to freeze out. The abundances of less volatile compounds such as ammonia, water, and hydrogen sulfide in the deep atmosphere are poorly known. They are probably also higher than solar values. Along with methane, trace amounts of various hydrocarbons are found in the stratosphere of Uranus, which are thought to be produced from methane by photolysis induced by the solar ultraviolet (UV) radiation. They include ethane (), acetylene (), methylacetylene (), and diacetylene (). Spectroscopy has also uncovered traces of water vapour, carbon monoxide, and carbon dioxide in the upper atmosphere, which can only originate from an external source such as infalling dust and comets.
Troposphere
The troposphere is the lowest and densest part of the atmosphere and is characterised by a decrease in temperature with altitude. The temperature falls from about at the base of the nominal troposphere at −300 km to at 50 km. The temperatures in the coldest upper region of the troposphere (the tropopause) actually vary in the range between depending on planetary latitude. The tropopause region is responsible for the vast majority of Uranus's thermal far infrared emissions, thus determining its effective temperature of .
The troposphere is thought to have a highly complex cloud structure; water clouds are hypothesised to lie in the pressure range of , ammonium hydrosulfide clouds in the range of , ammonia or hydrogen sulfide clouds at between and finally directly detected thin methane clouds at . The troposphere is a dynamic part of the atmosphere, exhibiting strong winds, bright clouds, and seasonal changes.
Upper atmosphere
The middle layer of the Uranian atmosphere is the stratosphere, where temperature generally increases with altitude from in the tropopause to between at the base of the thermosphere. The heating of the stratosphere is caused by absorption of solar UV and IR radiation by methane and other hydrocarbons, which form in this part of the atmosphere as a result of methane photolysis. Heat is also conducted from the hot thermosphere. The hydrocarbons occupy a relatively narrow layer at altitudes of between 100 and 300 km corresponding to a pressure range of 1,000 to 10 Pa and temperatures of between .
The most abundant hydrocarbons are methane, acetylene, and ethane with mixing ratios of around relative to hydrogen. The mixing ratio of carbon monoxide is similar at these altitudes. Heavier hydrocarbons and carbon dioxide have mixing ratios three orders of magnitude lower. The abundance ratio of water is around 7. Ethane and acetylene tend to condense in the colder lower part of the stratosphere and tropopause (below 10 mBar level) forming haze layers, which may be partly responsible for the bland appearance of Uranus. The concentration of hydrocarbons in the Uranian stratosphere above the haze is significantly lower than in the stratospheres of the other giant planets.
The outermost layer of the Uranian atmosphere is the thermosphere and corona, which has a uniform temperature of around to . The heat sources necessary to sustain such a high level are not understood, as neither the solar UV nor the auroral activity can provide the necessary energy to maintain these temperatures. The weak cooling efficiency due to the lack of hydrocarbons in the stratosphere above 0.1 mBar pressure levels may contribute too. In addition to molecular hydrogen, the thermosphere-corona contains many free hydrogen atoms. Their small mass and high temperatures explain why the corona extends as far as , or two Uranian radii, from its surface.
This extended corona is a unique feature of Uranus. Its effects include a drag on small particles orbiting Uranus, causing a general depletion of dust in the Uranian rings. The Uranian thermosphere, together with the upper part of the stratosphere, corresponds to the ionosphere of Uranus. Observations show that the ionosphere occupies altitudes from . The Uranian ionosphere is denser than that of either Saturn or Neptune, which may arise from the low concentration of hydrocarbons in the stratosphere. The ionosphere is mainly sustained by solar UV radiation and its density depends on the solar activity. Auroral activity is insignificant as compared to Jupiter and Saturn.
Climate
At ultraviolet and visible wavelengths, Uranus's atmosphere is bland in comparison to the other giant planets, even to Neptune, which it otherwise closely resembles. When Voyager 2 flew by Uranus in 1986, it observed a total of 10 cloud features across the entire planet. One proposed explanation for this dearth of features is that Uranus's internal heat is markedly lower than that of the other giant planets, being the coldest planet in the Solar System.
Banded structure, winds and clouds
In 1986, Voyager 2 found that the visible southern hemisphere of Uranus can be subdivided into two regions: a bright polar cap and dark equatorial bands. Their boundary is located at about −45° of latitude. A narrow band straddling the latitudinal range from −45 to −50° is the brightest large feature on its visible surface. It is called a southern "collar". The cap and collar are thought to be a dense region of methane clouds located within the pressure range of 1.3 to 2 bar. Besides the large-scale banded structure, Voyager 2 observed ten small bright clouds, most lying several degrees to the north from the collar. In all other respects, Uranus looked like a dynamically dead planet in 1986.
Voyager 2 arrived during the height of Uranus's southern summer and could not observe the northern hemisphere. At the beginning of the 21st century, when the northern polar region came into view, the Hubble Space Telescope (HST) and Keck telescope initially observed neither a collar nor a polar cap in the northern hemisphere. So Uranus appeared to be asymmetric: bright near the south pole and uniformly dark in the region north of the southern collar. In 2007, when Uranus passed its equinox, the southern collar almost disappeared, and a faint northern collar emerged near 45° of latitude. In 2023, a team employing the Very Large Array observed a dark collar at 80° latitude, and a bright spot at the north pole, indicating the presence of a polar vortex.
In the 1990s, the number of the observed bright cloud features grew considerably, partly because new high-resolution imaging techniques became available. Most were found in the northern hemisphere as it started to become visible. An early explanation—that bright clouds are easier to identify in its dark part, whereas in the southern hemisphere the bright collar masks them—was shown to be incorrect. Nevertheless, there are differences between the clouds of each hemisphere. The northern clouds are smaller, sharper and brighter. They appear to lie at a higher altitude. The lifetime of clouds spans several orders of magnitude. Some small clouds live for hours; at least one southern cloud may have persisted since the Voyager 2 flyby. Recent observation also discovered that cloud features on Uranus have a lot in common with those on Neptune. For example, the dark spots common on Neptune had never been observed on Uranus before 2006, when the first such feature dubbed Uranus Dark Spot was imaged. The speculation is that Uranus is becoming more Neptune-like during its equinoctial season.
The tracking of numerous cloud features allowed determination of zonal winds blowing in the upper troposphere of Uranus. At the equator winds are retrograde, which means that they blow in the reverse direction to the planetary rotation. Their speeds are from . Wind speeds increase with the distance from the equator, reaching zero values near ±20° latitude, where the troposphere's temperature minimum is located. Closer to the poles, the winds shift to a prograde direction, flowing with Uranus's rotation. Wind speeds continue to increase reaching maxima at ±60° latitude before falling to zero at the poles. Wind speeds at −40° latitude range from . Because the collar obscures all clouds below that parallel, speeds between it and the southern pole are impossible to measure. In contrast, in the northern hemisphere maximum speeds as high as are observed near +50° latitude.
In 1986, the Voyager 2 Planetary Radio Astronomy (PRA) experiment observed 140 lightning flashes, or Uranian electrostatic discharges with a frequency of 0.9-40 MHz. The UEDs were detected from 600,000 km of Uranus over 24 hours, most of which were not visible . However, microphysical modelling suggests that Uranian lightning occurs in convective storms occurring in deep troposphere water clouds. If this is the case, lightning will not be visible due to the thick cloud layers above the troposphere. The UEDs were detected from 600,000 km of Uranus, most of which were not visible . Uranian lightning has a power of around 108 W, emits 1×10^7 J - 2×10^7 J of energy, and lasts an average of 120 ms. There is a possibility that the power of Uranian lightning varies greatly with the seasons caused by changes in convection rates in the clouds The UEDs were detected from 600,000 km of Uranus, most of which were not visible. Uranian lightning is much more powerful than lightning on Earth and comparable to Jovian lightning. During the Ice Giant flybys, "Voyager 2" detected lightning more clearly on Uranus than on Neptune due to the planet's lower gravity and possible warmer deep atmosphere.
Seasonal variation
For a short period from March to May 2004, large clouds appeared in the Uranian atmosphere, giving it a Neptune-like appearance. Observations included record-breaking wind speeds of and a persistent thunderstorm referred to as "Fourth of July fireworks". On 23 August 2006, researchers at the Space Science Institute (Boulder, Colorado) and the University of Wisconsin observed a dark spot on Uranus's surface, giving scientists more insight into Uranus atmospheric activity. Why this sudden upsurge in activity occurred is not fully known, but it appears that Uranus's extreme axial tilt results in extreme seasonal variations in its weather. Determining the nature of this seasonal variation is difficult because good data on Uranus's atmosphere has existed for less than 84 years, or one full Uranian year. Photometry over the course of half a Uranian year (beginning in the 1950s) has shown regular variation in the brightness in two spectral bands, with maxima occurring at the solstices and minima occurring at the equinoxes. A similar periodic variation, with maxima at the solstices, has been noted in microwave measurements of the deep troposphere begun in the 1960s. Stratospheric temperature measurements beginning in the 1970s also showed maximum values near the 1986 solstice. The majority of this variability is thought to occur owing to changes in viewing geometry.
There are some indications that physical seasonal changes are happening in Uranus. Although Uranus is known to have a bright south polar region, the north pole is fairly dim, which is incompatible with the model of the seasonal change outlined above. During its previous northern solstice in 1944, Uranus displayed elevated levels of brightness, which suggests that the north pole was not always so dim. This information implies that the visible pole brightens some time before the solstice and darkens after the equinox. Detailed analysis of the visible and microwave data revealed that the periodical changes in brightness are not completely symmetrical around the solstices, which also indicates a change in the meridional albedo patterns.
In the 1990s, as Uranus moved away from its solstice, Hubble and ground-based telescopes revealed that the south polar cap darkened noticeably (except the southern collar, which remained bright), whereas the northern hemisphere demonstrated increasing activity, such as cloud formations and stronger winds, bolstering expectations that it should brighten soon. This indeed happened in 2007 when it passed an equinox: a faint northern polar collar arose, and the southern collar became nearly invisible, although the zonal wind profile remained slightly asymmetric, with northern winds being somewhat slower than southern.
The mechanism of these physical changes is still not clear. Near the summer and winter solstices, Uranus's hemispheres lie alternately either in full glare of the Sun's rays or facing deep space. The brightening of the sunlit hemisphere is thought to result from the local thickening of the methane clouds and haze layers located in the troposphere. The bright collar at −45° latitude is also connected with methane clouds. Other changes in the southern polar region can be explained by changes in the lower cloud layers. The variation of the microwave emission from Uranus is probably caused by changes in the deep tropospheric circulation, because thick polar clouds and haze may inhibit convection. Now that the spring and autumn equinoxes are arriving on Uranus, the dynamics are changing and convection can occur again.
Magnetosphere
Before the arrival of Voyager 2, no measurements of the Uranian magnetosphere had been taken, so its nature remained a mystery. Before 1986, scientists had expected the magnetic field of Uranus to be in line with the solar wind, because it would then align with Uranus's poles that lie in the ecliptic.
Voyagers observations revealed that Uranus's magnetic field is peculiar, both because it does not originate from its geometric centre, and because it is tilted at 59° from the axis of rotation. In fact, the magnetic dipole is shifted from Uranus's centre towards the south rotational pole by as much as one-third of the planetary radius. This unusual geometry results in a highly asymmetric magnetosphere, where the magnetic field strength on the surface in the southern hemisphere can be as low as 0.1 gauss (10 μT), whereas in the northern hemisphere it can be as high as 1.1 gauss (110 μT). The average field at the surface is 0.23 gauss (23 μT).
Studies of Voyager 2 data in 2017 suggest that this asymmetry causes Uranus's magnetosphere to connect with the solar wind once a Uranian day, opening the planet to the Sun's particles. In comparison, the magnetic field of Earth is roughly as strong at either pole, and its "magnetic equator" is roughly parallel with its geographical equator. The dipole moment of Uranus is 50 times that of Earth. Neptune has a similarly displaced and tilted magnetic field, suggesting that this may be a common feature of ice giants. One hypothesis is that, unlike the magnetic fields of the terrestrial and gas giants, which are generated within their cores, the ice giants' magnetic fields are generated by motion at relatively shallow depths, for instance, in the water–ammonia ocean. Another possible explanation for the magnetosphere's alignment is that there are oceans of liquid diamond in Uranus's interior that would deter the magnetic field.
It is, however, unclear whether the observed asymmetry of Uranus' magnetic field represents the typical state of the magnetosphere, or a coincidence of observing it during unusual space weather conditions. A post-analysis of Voyager data from 2024 suggests that the strongly asymmetric shape of the magnetosphere observed during the fly-by represents an anomalous state, as the measured values of solar wind density at the time were unusually high, which could have compressed Uranus' magnetosphere. The interaction with the solar wind event could also explain the apparent paradox of presence of strong electron radiation belts despite the otherwise low magnetospheric plasma density measured. Such conditions are estimated to occur less than 5% of the time.
Despite its curious alignment, in other respects the Uranian magnetosphere is like those of other planets: it has a bow shock at about 23 Uranian radii ahead of it, a magnetopause at 18 Uranian radii, a fully developed magnetotail, and radiation belts. Overall, the structure of Uranus's magnetosphere is different from Jupiter's and more similar to Saturn's. Uranus's magnetotail trails behind it into space for millions of kilometres and is twisted by its sideways rotation into a long corkscrew.Uranus's magnetosphere contains charged particles: mainly protons and electrons, with a small amount of H2+ ions. Many of these particles probably derive from the thermosphere. The ion and electron energies can be as high as 4 and 1.2 megaelectronvolts, respectively. The density of low-energy (below 1 kiloelectronvolt) ions in the inner magnetosphere is about 2 cm−3. The particle population is strongly affected by the Uranian moons, which sweep through the magnetosphere, leaving noticeable gaps. The particle flux is high enough to cause darkening or space weathering of their surfaces on an astronomically rapid timescale of 100,000 years. This may be the cause of the uniformly dark colouration of the Uranian satellites and rings.
Uranus has relatively well developed aurorae, which are seen as bright arcs around both magnetic poles. Unlike Jupiter's, Uranus's aurorae seem to be insignificant for the energy balance of the planetary thermosphere. They, or rather their trihydrogen cations' infrared spectral emissions, have been studied in-depth as of late 2023.
In March 2020, NASA astronomers reported the detection of a large atmospheric magnetic bubble, also known as a plasmoid, released into outer space from the planet Uranus, after reevaluating old data recorded by the Voyager 2 space probe during a flyby of the planet in 1986.
Moons
Uranus has 28 known natural satellites. The names of these satellites are chosen from characters in the works of Shakespeare and Alexander Pope. The five main satellites are Miranda, Ariel, Umbriel, Titania, and Oberon. The Uranian satellite system is the least massive among those of the giant planets; the combined mass of the five major satellites would be less than half that of Triton (largest moon of Neptune) alone. The largest of Uranus's satellites, Titania, has a radius of only , or less than half that of the Moon, but slightly more than Rhea, the second-largest satellite of Saturn, making Titania the eighth-largest moon in the Solar System. Uranus's satellites have relatively low albedos; ranging from 0.20 for Umbriel to 0.35 for Ariel (in green light). They are ice–rock conglomerates composed of roughly 50% ice and 50% rock. The ice may include ammonia and carbon dioxide.
Among the Uranian satellites, Ariel appears to have the youngest surface, with the fewest impact craters, and Umbriel the oldest. Miranda has fault canyons deep, terraced layers, and a chaotic variation in surface ages and features. Miranda's past geologic activity is thought to have been driven by tidal heating at a time when its orbit was more eccentric than currently, probably as a result of a former 3:1 orbital resonance with Umbriel. Extensional processes associated with upwelling diapirs are the likely origin of Miranda's 'racetrack'-like coronae. Ariel is thought to have once been held in a 4:1 resonance with Titania.
Uranus has at least one horseshoe orbiter occupying the Sun–Uranus Lagrangian point—a gravitationally unstable region at 180° in its orbit, 83982 Crantor. Crantor moves inside Uranus's co-orbital region on a complex, temporary horseshoe orbit. is also a promising Uranus horseshoe librator candidate.
Rings
The Uranian rings are composed of extremely dark particles, which vary in size from micrometres to a fraction of a metre. Thirteen distinct rings are presently known, the brightest being the ε ring. All except the two rings of Uranus are extremely narrow—they are usually a few kilometres wide. The rings are probably quite young; the dynamics considerations indicate that they did not form with Uranus. The matter in the rings may once have been part of a moon (or moons) that was shattered by high-speed impacts. From numerous pieces of debris that formed as a result of those impacts, only a few particles survived, in stable zones corresponding to the locations of the present rings.
William Herschel described a possible ring around Uranus in 1789. This sighting is generally considered doubtful, because the rings are quite faint, and in the two following centuries none were noted by other observers. Still, Herschel made an accurate description of the epsilon ring's size, its angle relative to Earth, its red colour, and its apparent changes as Uranus travelled around the Sun. The ring system was definitively discovered on 10 March 1977 by James L. Elliot, Edward W. Dunham, and Jessica Mink using the Kuiper Airborne Observatory. The discovery was serendipitous; they planned to use the occultation of the star SAO 158687 (also known as HD 128598) by Uranus to study its atmosphere. When their observations were analysed, they found that the star had disappeared briefly from view five times both before and after it disappeared behind Uranus. They concluded that there must be a ring system around Uranus. Later, they detected four additional rings. The rings were directly imaged when Voyager 2 passed Uranus in 1986. Voyager 2 also discovered two additional faint rings, bringing the total number to eleven.
In December 2005, the Hubble Space Telescope detected a pair of previously unknown rings. The largest is located twice as far from Uranus as the previously known rings. These new rings are so far from Uranus that they are called the "outer" ring system. Hubble also spotted two small satellites, one of which, Mab, shares its orbit with the outermost newly discovered ring. The new rings bring the total number of Uranian rings to 13. In April 2006, images of the new rings from the Keck Observatory yielded the colours of the outer rings: the outermost is blue and the other one red. One hypothesis concerning the outer ring's blue colour is that it is composed of minute particles of water ice from the surface of Mab that are small enough to scatter blue light. In contrast, Uranus's inner rings appear grey.
Although the Uranian rings are very difficult to directly observe from Earth, advances in digital imaging have allowed several amateur astronomers to successfully photograph the rings with red or infrared filters; telescopes with apertures as small as may be able to detect the rings with proper imaging equipment.
Exploration
Launched in 1977, Voyager 2 made its closest approach to Uranus on 24 January 1986, coming within of the cloudtops, before continuing its journey to Neptune. The spacecraft studied the structure and chemical composition of Uranus's atmosphere, including its unique weather, caused by its extreme axial tilt. It made the first detailed investigations of its five largest moons and discovered 10 new ones. Voyager 2 examined all nine of the system's known rings and discovered two more. It also studied the magnetic field, its irregular structure, its tilt and its unique corkscrew magnetotail caused by Uranus's sideways orientation.
No other spacecraft has flown by Uranus since then, though there have been many proposed missions to revisit the Uranus system. The possibility of sending the Cassini spacecraft from Saturn to Uranus was evaluated during a mission extension planning phase in 2009, but was ultimately rejected in favour of destroying it in the Saturnian atmosphere, as it would have taken about twenty years to get to the Uranian system after departing Saturn. A Uranus entry probe could use Pioneer Venus Multiprobe heritage and descend to 1–5 atmospheres. A Uranus orbiter and probe was recommended by the 2013–2022 Planetary Science Decadal Survey published in 2011; the proposal envisaged launch during 2020–2023 and a 13-year cruise to Uranus. The committee's opinion was reaffirmed in 2022, when a Uranus probe/orbiter mission was placed at the highest priority, due to the lack of knowledge about ice giants. Most recently, the CNSA's Tianwen-4 Jupiter orbiter, launching in 2029, is planned to have a subprobe that will detach and get a gravity assist instead of entering orbit, flying by Uranus in March 2045 before heading to interstellar space. China also has plans for a potential Tianwen-5 that may orbit either Uranus or Neptune that have yet to come to fruition.
In culture
In modern astrology, the planet Uranus (symbol ) is the ruling planet of Aquarius; prior to the discovery of Uranus, the ruling planet of Aquarius was Saturn. Because Uranus is cyan and Uranus is associated with electricity, the colour electric blue, which is close to cyan, is associated with the sign Aquarius.
The chemical element uranium, discovered in 1789 by the German chemist Martin Heinrich Klaproth, was named after the then-newly discovered Uranus.
Lydia Sigourney included her poem in her 1827 collection of poetry.
"Uranus, the Magician" is a movement in Gustav Holst's orchestral suite The Planets, written between 1914 and 1916.
Operation Uranus was the successful military operation in World War II by the Red Army to take back Stalingrad and marked the turning point in the land war against the Wehrmacht.
The lines "Then felt I like some watcher of the skies/When a new planet swims into his ken", from John Keats's "On First Looking into Chapman's Homer", are a reference to Herschel's discovery of Uranus.
| Physical sciences | Astronomy | null |
44481 | https://en.wikipedia.org/wiki/Tuff | Tuff | Tuff is a type of rock made of volcanic ash ejected from a vent during a volcanic eruption. Following ejection and deposition, the ash is lithified into a solid rock. Rock that contains greater than 75% ash is considered tuff, while rock containing 25% to 75% ash is described as tuffaceous (for example, tuffaceous sandstone). A pyroclastic rock containing 25–75% volcanic bombs and/or volcanic blocks is called tuff breccia. Tuff composed of sandy volcanic material can be referred to as volcanic sandstone.
Tuff is a relatively soft rock, so it has been used for construction since ancient times. Because it is common in Italy, the Romans used it often for construction. The Rapa Nui people used it to make most of the moai statues on Easter Island.
Tuff can be classified as either igneous or sedimentary rock. It is usually studied in the context of igneous petrology, although it is sometimes described using sedimentological terms.
Tuff is often erroneously called tufa in guidebooks and in television programs but tufa is a form of travertine.
Volcanic ash
The material that is expelled in a volcanic eruption can be classified into three types:
Volcanic gases, a mixture made mostly of steam, carbon dioxide, and a sulfur compound (either sulfur dioxide, SO2, or hydrogen sulfide, H2S, depending on the temperature)
Lava, the name of magma when it emerges and flows over the surface
Tephra, particles of solid material of all shapes and sizes ejected and thrown through the air
Tephra is made when magma inside the volcano is blown apart by the rapid expansion of hot volcanic gases. Magma commonly explodes as the gas dissolved in it comes out of solution as the pressure decreases when it flows to the surface. These violent explosions produce particles of material that can then fly from the volcano. Solid particles smaller than 2 mm in diameter (sand-sized or smaller) are called volcanic ash.
Volcanic ash is further divided into fine ash, with particle sizes smaller than 0.0625 mm in diameter, and coarse ash, with particle sizes between 0.0625 mm and 2 mm in diameter. Tuff is correspondingly divided into coarse tuff (coarse ash tuff) and fine tuff (fine ash tuff or dust tuff). Consolidated tephra composed mostly of coarser particles is called lapillistone (particles 2 mm to 64 mm in diameter) or agglomerate or pyroclastic breccia (particles over 64 mm in diameter) rather than tuff.
Volcanic ash can vary greatly in composition, and so tuffs are further classified by the composition of the ash from which they formed. Ash from high-silica volcanism, particularly in ash flows, consists mainly of shards of volcanic glass, and tuff formed predominantly from glass shards is described as vitric tuff. The glass shards are typically either irregular in shape or are roughly triangular with convex sides. They are the shattered walls of countless small bubbles that formed in the magma as dissolved gases rapidly came out of solution.
Tuffs formed from ash consisting predominantly of individual crystals are described as crystal tuffs, while those formed from ash consisting predominantly of pulverized rock fragments are described as lithic tuffs.
The chemical composition of volcanic ash reflects the entire range of volcanic rock chemistry, from high-silica rhyolitic ash to low-silica basaltic ash, and tuffs are likewise described as rhyolitic, andesitic, basaltic, and so on.
Transport and lithification
The most straightforward way for volcanic ash to move away from the vent is as ash clouds that are part of an eruption column. These fall to the surface as fallout deposits that are characteristically well-sorted and tend to form a blanket of uniform thickness across terrain. Column collapse results in a more spectacular and destructive form of transport, which takes the form of pyroclastic flows and surges that characteristically are poorly sorted and pool in low terrain. Surge deposits sometimes show sedimentary structures typical of high-velocity flow, such as dunes and antidunes. Volcanic ash already deposited on the surface can be transported as mud flows (lahars) when mingled with water from rainfall or through eruption into a body of water or ice.
Particles of volcanic ash that are sufficiently hot will weld together after settling to the surface, producing a welded tuff. Welding requires temperatures in excess of . If the rock contains scattered, pea-sized fragments or fiamme in it, it is called a welded lapilli tuff. Welded tuffs (and welded lapilli tuffs) can be of fallout origin, or deposited from ash flows, as in the case of ignimbrites. During welding, the glass shards and pumice fragments adhere together (necking at point contacts), deform, and compact together, resulting in a eutaxitic fabric. Welded tuff is commonly rhyolitic in composition, but examples of all compositions are known.
A sequence of ash flows may consist of multiple cooling units. These can be distinguished by the degree of welding. The base of a cooling unit is typically unwelded due to chilling from the underlying cold surface, and the degree of welding and of secondary reactions from fluids in the flow increases upwards towards the center of the flow. Welding decreases towards the top of the cooling unit, where the unit cools more rapidly. The intensity of welding may also decrease towards areas in which the deposit is thinner, and with distance from source.
Cooler pyroclastic flows are unwelded and the ash sheets deposited by them are relatively unconsolidated. However, cooled volcanic ash can quickly become lithified because it usually has a high content of volcanic glass. This is a thermodynamically unstable material that reacts rapidly with ground water or sea water, which leaches alkali metals and calcium from the glass. New minerals, such as zeolites, clays, and calcite, crystallize from the dissolved substances and cement the tuff.
Tuffs are further classified by their depositional environment, such as lacustrine tuff, subaerial tuff, or submarine tuff, or by the mechanism by which the ash was transported, such as fallout tuff or ash flow tuff. Reworked tuffs, formed by erosion and redeposition of ash deposits, are usually described by the transport agent, such as aeolian tuff or fluvial tuff.
Occurrences
Tuffs have the potential to be deposited wherever explosive volcanism takes place, and so have a wide distribution in location and age.
High-silica volcanism
Rhyolite tuffs contain pumiceous, glassy fragments and small scoriae with quartz, alkali feldspar, biotite, etc. Iceland, Lipari, Hungary, the Basin and Range of the American southwest, and New Zealand are among the areas where such tuffs are prominent. In the ancient rocks of Wales, Charnwood, etc., similar tuffs are known, but in all cases, they are greatly changed by silicification (which has filled them with opal, chalcedony, and quartz) and by devitrification. The frequent presence of rounded corroded quartz crystals, such as occur in rhyolitic lavas, helps to demonstrate their real nature.
Welded ignimbrites can be highly voluminous, such as the Lava Creek Tuff erupted from Yellowstone Caldera in Wyoming 631,000 years ago. This tuff had an original volume of at least . Lava Creek tuff is known to be at least 1000 times as large as the deposits of the 1980 eruption of Mount St. Helens, and it had a Volcanic Explosivity Index (VEI) of 8, greater than any eruption known in the last 10,000 years. Ash flow tuffs cover of the North Island of New Zealand and about of Nevada. Ash flow tuffs are the only volcanic product with volumes rivaling those of flood basalts.
The Tioga Bentonite of the northeastern United States varies in composition from crystal tuff to tuffaceous shale. It was deposited as ash carried by wind that fell out over the sea and settled to the bottom. It is Devonian in age and likely came from a vent in central Virginia, where the tuff reaches its maximum thickness of about .
Alkaline volcanism
Trachyte tuffs contain little or no quartz, but much sanidine or anorthoclase and sometimes oligoclase feldspar, with occasional biotite, augite, and hornblende. In weathering, they often change to soft red or yellow claystones, rich in kaolin with secondary quartz. Recent trachyte tuffs are found on the Rhine (at Siebengebirge), in Ischia and near Naples. Trachyte-carbonatite tuffs have been identified in the East African Rift. Alkaline crystal tuffs have been reported from Rio de Janeiro.
Intermediate volcanism
Andesitic tuffs are exceedingly common. They occur along the whole chain of the Cordilleras and Andes, in the West Indies, New Zealand, Japan, etc. In the Lake District, North Wales, Lorne, the Pentland Hills, the Cheviots, and many other districts of Great Britain, ancient rocks of exactly similar nature are abundant. In color, they are red or brown; their scoriae fragments are of all sizes from huge blocks down to minute granular dust. The cavities are filled with many secondary minerals, such as calcite, chlorite, quartz, epidote, or chalcedony; in microscopic sections, though, the nature of the original lava can nearly always be made out from the shapes and properties of the little crystals which occur in the decomposed glassy base. Even in the smallest details, these ancient tuffs have a complete resemblance to the modern ash beds of Cotopaxi, Krakatoa, and Mont Pelé.
Mafic volcanism
Mafic volcanism typically takes the form of Hawaiian eruptions that are nonexplosive and produce little ash. However, interaction between basaltic magma and groundwater or sea water results in hydromagmatic explosions that produce abundant ash. These deposit ash cones that subsequently can become cemented into tuff cones. Diamond Head, Hawaii, is an example of a tuff cone, as is the island of Ka'ula. The glassy basaltic ash produced in such eruptions rapidly alters to palagonite as part of the process of lithification.
Although conventional mafic volcanism produce little ash, such ash as is formed may accumulate locally as significant deposits. An example is the Pahala ash of Hawaii island, which locally is as thick as . These deposits also rapidly alter to palagonite, and eventually weather to laterite.
Basaltic tuffs are also found in Skye, Mull, Antrim, and other places, where Paleogene volcanic rocks are found; in Scotland, Derbyshire, and Ireland among the Carboniferous strata, and among the still older rocks of the Lake District, the southern uplands of Scotland, and Wales. They are black, dark green, or red in colour; vary greatly in coarseness, some being full of round spongy bombs a foot or more in diameter; and being often submarine, may contain shale, sandstone, grit, and other sedimentary material, and are occasionally fossiliferous. Recent basaltic tuffs are found in Iceland, the Faroe Islands, Jan Mayen, Sicily, the Hawaiian Islands, Samoa, etc. When weathered, they are filled with calcite, chlorite, serpentine, and especially where the lavas contain nepheline or leucite, are often rich in zeolites, such as analcite, prehnite, natrolite, scolecite, chabazite, heulandite, etc.
Ultramafic volcanism
Ultramafic tuffs are extremely rare; their characteristic is the abundance of olivine or serpentine and the scarcity or absence of feldspar and quartz.
Kimberlites
Occurrences of ultramafic tuff include surface deposits of kimberlite at maars in the diamond-fields of southern Africa and other regions. The principal variety of kimberlite is a dark bluish-green, serpentine-rich breccia (blue-ground) which, when thoroughly oxidized and weathered, becomes a friable brown or yellow mass (the "yellow-ground"). These breccias were emplaced as gas–solid mixtures and are typically preserved and mined in diatremes that form intrusive pipe-like structures. At depth, some kimberlite breccias grade into root zones of dikes made of unfragmented rock. At the surface, ultramafic tuffs may occur in maar deposits. Because kimberlites are the most common igneous source of diamonds, the transitions from maar to diatreme to root-zone dikes have been studied in detail. Diatreme-facies kimberlite is more properly called an ultramafic breccia rather than a tuff.
Komatiites
Komatiite tuffs are found, for example, in the greenstone belts of Canada and South Africa.
Folding and metamorphism
In course of time, changes other than weathering may overtake tuff deposits. Sometimes, they are involved in folding and become sheared and cleaved. Many of the green slates of the English Lake District are finely cleaved ashes. In Charnwood Forest also, the tuffs are slaty and cleaved. The green color is due to the large development of chlorite. Among the crystalline schists of many regions, green beds or green schists occur, which consist of quartz, hornblende, chlorite or biotite, iron oxides, feldspar, etc., and are probably recrystallized or metamorphosed tuffs. They often accompany masses of epidiorite and hornblende – schists which are the corresponding lavas and sills. Some chlorite-schists also are probably altered beds of volcanic tuff. The "Schalsteins" of Devon and Germany include many cleaved and partly recrystallized ash-beds, some of which still retain their fragmental structure, though their lapilli are flattened and drawn out. Their steam cavities are usually filled with calcite, but sometimes with quartz. The more completely altered forms of these rocks are platy, green chloritic schists; in these, however, structures indicating their original volcanic nature only sparingly occur. These are intermediate stages between cleaved tuffs and crystalline schists.
Importance
The primary economic value of tuff is as a building material. In the ancient world, tuff's relative softness meant that it was commonly used for construction where it was available.
Italy
Tuff is common in Italy, and the Romans used it for many buildings and bridges. For example, the whole port of the island of Ventotene (still in use), was carved from tuff. The Servian Wall, built to defend the city of Rome in the fourth century BC, is also built almost entirely from tuff. The Romans also cut tuff into small, rectangular stones that they used to create walls in a pattern known as opus reticulatum.
Peperino has been used in Rome and Naples as a building stone, is a trachyte tuff. Pozzolana also is a decomposed tuff, but of basic character, originally obtained near Naples and used as a cement, but this name is now applied to a number of substances not always of identical character. In the historical architecture of Naples, Neapolitan yellow tuff is the most used building material. Piperno ignimbrite tuff was also used widely in Naples and Campania.
Germany
In the Eifel region of Germany, a trachytic, pumiceous tuff called trass has been extensively worked as a hydraulic mortar. Tuff of the Eifel region of Germany has been widely used for construction of railroad stations and other buildings in Frankfurt, Hamburg, and other large cities. Construction using the Rochlitz Porphyr, can be seen in the Mannerist-style sculpted portal outside the chapel entrance in Colditz Castle. The trade name Rochlitz Porphyr is the traditional designation for a dimension stone of Saxony with an architectural history over 1,000 years in Germany. The quarries are located near Rochlitz.
United States
Yucca Mountain nuclear waste repository, a U.S. Department of Energy terminal storage facility for spent nuclear reactor and other radioactive waste, is in tuff and ignimbrite in the Basin and Range Province in Nevada. In Napa Valley and Sonoma Valley, California, areas made of tuff are routinely excavated for storage of wine barrels.
Rapa Nui
Tuff from Rano Raraku was used by the Rapa Nui people of Easter Island to make the vast majority of their famous moai statues.
Armenia
Tuff is used extensively in Armenia and Armenian architecture. It is the dominant type of stone used in construction in Armenia's capital Yerevan, Gyumri, Armenia's second largest city, and Ani, the country's medieval capital, now in Turkey. A small village in Armenia was renamed Tufashen (literally "village of tuff") in 1946.
Tephrochronology
Tuffs are deposited geologically instantaneously and often over a large region. This makes them highly useful as time-stratigraphic markers. The use of tuffs and other tephra deposits in this manner is known as tephrochronology and is particularly useful for Quaternary chronostratigraphy. Individual tuff beds can be "fingerprinted" by their chemical composition and phenocryst assemblages. Absolute ages for tuff beds can be determined by K-Ar, Ar-Ar, or carbon-14 dating. Zircon grains found in many tuffs are highly durable and can survive even metamorphism of the host tuff to schist, allowing absolute ages to be assigned to ancient metamorphic rocks. For example, dating of zircons in a metamorphosed tuff bed in the Pilar Formation provided some of the first evidence for the Picuris orogeny.
Etymology
The word tuff is derived from the Italian tufo.
| Physical sciences | Petrology | null |
44495 | https://en.wikipedia.org/wiki/Linear%20motor | Linear motor | A linear motor is an electric motor that has had its stator and rotor "unrolled", thus, instead of producing a torque (rotation), it produces a linear force along its length. However, linear motors are not necessarily straight. Characteristically, a linear motor's active section has ends, whereas more conventional motors are arranged as a continuous loop.
A typical mode of operation is as a Lorentz-type actuator, in which the applied force is linearly proportional to the current and the magnetic field .
Linear motors are most commonly found in high accuracy engineering applications.
Many designs have been put forward for linear motors, falling into two major categories, low-acceleration and high-acceleration linear motors. Low-acceleration linear motors are suitable for maglev trains and other ground-based transportation applications. High-acceleration linear motors are normally rather short, and are designed to accelerate an object to a very high speed; for example, see the coilgun.
High-acceleration linear motors are typically used in studies of hypervelocity collisions, as weapons, or as mass drivers for spacecraft propulsion. They are usually of the AC linear induction motor (LIM) design with an active three-phase winding on one side of the air-gap and a passive conductor plate on the other side. However, the direct current homopolar linear motor railgun is another high acceleration linear motor design. The low-acceleration, high speed and high power motors are usually of the linear synchronous motor (LSM) design, with an active winding on one side of the air-gap and an array of alternate-pole magnets on the other side. These magnets can be permanent magnets or electromagnets. The motor for the Shanghai maglev train, for instance, is an LSM.
Types
Brushless
Brushless linear motors are members of the Synchronous motor family. They are typically used in standard linear stages or integrated into custom, high performance positioning systems. Invented in the late 1980s by Anwar Chitayat at Anorad Corporation, now Rockwell Automation, and helped improve the throughput and quality of industrial manufacturing processes.
Brush
Brushed linear motors were used in industrial automation applications prior to the invention of Brushless linear motors. Compared with three phase brushless motors, which are typically being used today, brush motors operate on a single phase. Brush linear motors have a lower cost since they do not need moving cables or three phase servo drives. However, they require higher maintenance since their brushes wear out.
Synchronous
In this design the rate of movement of the magnetic field is controlled, usually electronically, to track the motion of the rotor. For cost reasons synchronous linear motors rarely use commutators, so the rotor often contains permanent magnets, or soft iron. Examples include coilguns and the motors used on some maglev systems, as well as many other linear motors. In high precision industrial automation linear motors are typically configured with a magnet stator and a moving coil. A Hall effect sensor is attached to the rotor to track the magnetic flux of the stator. The electric current is typically provided from a stationary servo drive to the moving coil by a moving cable inside a cable carrier.
Induction
In this design, the force is produced by a moving linear magnetic field acting on conductors in the field. Any conductor, be it a loop, a coil or simply a piece of plate metal, that is placed in this field will have eddy currents induced in it thus creating an opposing magnetic field, in accordance with Lenz's law. The two opposing fields will repel each other, thus creating motion as the magnetic field sweeps through the metal.
Homopolar
In this design a large current is passed through a metal sabot across sliding contacts that are fed by two rails. The magnetic field this generates causes the metal to be projected along the rails.
Tubular
Efficient and compact design applicable to the replacement of pneumatic cylinders.
Piezoelectric
Piezoelectric drive is often used to drive small linear motors.
History
Low acceleration
The history of linear electric motors can be traced back at least as far as the 1840s, to the work of Charles Wheatstone at King's College London, but Wheatstone's model was too inefficient to be practical. A feasible linear induction motor is described in (1905 - inventor Alfred Zehden of Frankfurt-am-Main), for driving trains or lifts. The German engineer Hermann Kemper built a working model in 1935. In the late 1940s, Dr. Eric Laithwaite of Manchester University, later Professor of Heavy Electrical Engineering at Imperial College in London developed the first full-size working model.
In a single sided version the magnetic repulsion forces the conductor away from the stator, levitating it, and carrying it along in the direction of the moving magnetic field. He called the later versions of it magnetic river. The technologies would later be applied, in the 1984, Air-Rail Link shuttle, between Birmingham's airport and an adjacent train station.
Because of these properties, linear motors are often used in maglev propulsion, as in the Japanese Linimo magnetic levitation train line near Nagoya. However, linear motors have been used independently of magnetic levitation, as in the Bombardier Innovia Metro systems worldwide and a number of modern Japanese subways, including Tokyo's Toei Ōedo Line.
Similar technology is also used in some roller coasters with modifications but, at present, is still impractical on street running trams, although this, in theory, could be done by burying it in a slotted conduit.
Outside of public transportation, vertical linear motors have been proposed as lifting mechanisms in deep mines, and the use of linear motors is growing in motion control applications. They are also often used on sliding doors, such as those of low floor trams such as the Alstom Citadis and the Socimi Eurotram. Dual axis linear motors also exist. These specialized devices have been used to provide direct X-Y motion for precision laser cutting of cloth and sheet metal, automated drafting, and cable forming. Most linear motors in use are LIM (linear induction motor), or LSM (linear synchronous motor). Linear DC motors are not used due to their higher cost and linear SRM suffers from poor thrust. So for long runs in traction LIM is mostly preferred and for short runs LSM is mostly preferred.
High acceleration
High-acceleration linear motors have been suggested for a number of uses.
They have been considered for use as weapons, since current armour-piercing ammunition tends to consist of small rounds with very high kinetic energy, for which just such motors are suitable. Many amusement park launched roller coasters now use linear induction motors to propel the train at a high speed, as an alternative to using a lift hill.
The United States Navy is also using linear induction motors in the Electromagnetic Aircraft Launch System that will replace traditional steam catapults on future aircraft carriers. They have also been suggested for use in spacecraft propulsion. In this context they are usually called mass drivers. The simplest way to use mass drivers for spacecraft propulsion would be to build a large mass driver that can accelerate cargo up to escape velocity, though RLV launch assist like StarTram to low Earth orbit has also been investigated.
High-acceleration linear motors are difficult to design for a number of reasons. They require large amounts of energy in very short periods of time. One rocket launcher design calls for 300 GJ for each launch in the space of less than a second. Normal electrical generators are not designed for this kind of load, but short-term electrical energy storage methods can be used. Capacitors are bulky and expensive but can supply large amounts of energy quickly. Homopolar generators can be used to convert the kinetic energy of a flywheel into electric energy very rapidly. High-acceleration linear motors also require very strong magnetic fields; in fact, the magnetic fields are often too strong to permit the use of superconductors. However, with careful design, this need not be a major problem.
Two different basic designs have been invented for high-acceleration linear motors: railguns and coilguns.
Usage
Linear motors are commonly used for actuating high performance industrial automation equipment. Their advantage, unlike any other commonly used actuator, such as a ball screw, timing belt, or rack and pinion, is that they provide any combination of high precision, high velocity, high force and long travel.
Linear motors are widely used. One of the major uses of linear motors is for propelling the shuttle in looms.
A linear motor has been used for sliding doors and various similar actuators. They have been used for baggage handling and even large-scale bulk materials transport.
Linear motors are sometimes used to create rotary motion. For example, they have been used at observatories to deal with the large radius of curvature.
Linear motors may also be used as an alternative to conventional chain-run lift hills for roller coasters. The coaster Maverick at Cedar Point uses one such linear motor in place of a chain lift.
A linear motor has been used to accelerate cars for crash tests.
Industrial automation
The combination of high precision, high velocity, high force, and long travel makes brushless linear motors attractive for driving industrial automations equipment. They serve industries and applications such as semiconductor steppers, electronics surface-mount technology, automotive cartesian coordinate robots, aerospace chemical milling, optics electron microscope, healthcare laboratory automation, food and beverage pick and place.
Machine tools
Synchronous linear motor actuators, used in machine tools, provide high force, high velocity, high precision and high dynamic stiffness, resulting in high smoothness of motion and low settling time. They may reach velocities of 2 m/s and micron-level accuracies, with short cycle times and a smooth surface finish.
Train propulsion
Conventional rails
All of the following applications are in rapid transit and have the active part of the motor in the cars.
Bombardier Innovia Metro
Originally developed in the late 1970s by UTDC in Canada as the Intermediate Capacity Transit System (ICTS). A test track was constructed in Millhaven, Ontario, for extensive testing of prototype cars, after which three lines were constructed:
Line 3 Scarborough in Toronto (opened 1985; closed 2023)
Expo Line of the Vancouver SkyTrain (opened 1985 and extended in 1994)
Detroit People Mover in Detroit (opened 1987)
ICTS was sold to Bombardier Transportation in 1991 and later known as Advanced Rapid Transit (ART) before adopting its current branding in 2011. Since then, several more installations have been made:
Kelana Jaya Line in Kuala Lumpur (opened 1998 and extended in 2016)
Millennium Line of the Vancouver SkyTrain (opened 2002 and extended in 2016)
AirTrain JFK in New York (opened 2003)
Airport Express (Beijing Subway) (opened 2008)
Everline in Yongin, South Korea (opened 2013)
All Innovia Metro systems use third rail electrification.
Japanese Linear Metro
One of the biggest challenges faced by Japanese railway engineers in the 1970s to the 1980s was the ever increasing construction costs of subways. In response, the Japan Subway Association began studying on the feasibility of the "mini-metro" for meeting urban traffic demand in 1979. In 1981, the Japan Railway Engineering Association studied on the use of linear induction motors for such small-profile subways and by 1984 was investigating on the practical applications of linear motors for urban rail with the Japanese Ministry of Land, Infrastructure, Transport and Tourism. In 1988, a successful demonstration was made with the Limtrain at Saitama and influenced the eventual adoption of the linear motor for the Nagahori Tsurumi-ryokuchi Line in Osaka and Toei Line 12 (present-day Toei Oedo Line) in Tokyo.
To date, the following subway lines in Japan use linear motors and use overhead lines for power collection:
Two Osaka Metro lines in Osaka:
Nagahori Tsurumi-ryokuchi Line (opened 1990)
Imazatosuji Line (opened 2006)
Toei Ōedo Line in Tokyo (opened 2000)
Kaigan Line of the Kobe Municipal Subway (opened 2001)
Nanakuma Line of the Fukuoka City Subway (opened 2005)
Yokohama Municipal Subway Green Line (opened 2008)
Sendai Subway Tōzai Line (opened 2015)
In addition, Kawasaki Heavy Industries has also exported the Linear Metro to the Guangzhou Metro in China; all of the Linear Metro lines in Guangzhou use third rail electrification:
Line 4 (opened 2005)
Line 5 (opened 2009).
Line 6 (opened 2013)
Monorail
There is at least one known monorail system which is not magnetically levitated, but nonetheless uses linear motors. This is the Moscow Monorail. Originally, traditional motors and wheels were to be used. However, it was discovered during test runs that the proposed motors and wheels would fail to provide adequate traction under some conditions, for example, when ice appeared on the rail. Hence, wheels are still used, but the trains use linear motors to accelerate and slow down. This is possibly the only use of such a combination, due to the lack of such requirements for other train systems.
The TELMAGV is a prototype of a monorail system that is also not magnetically levitated but uses linear motors.
Magnetic levitation
High-speed trains:
Transrapid: first commercial use in Shanghai (opened in 2004)
SCMaglev, under construction in Japan (fastest train in the world, planned to open by 2027)
Rapid transit:
Birmingham Airport, UK (opened 1984, closed 1995)
M-Bahn in Berlin, Germany (opened in 1989, closed in 1991)
Daejeon EXPO, Korea (ran only 1993)
HSST: Linimo line in Aichi Prefecture, Japan (opened 2005)
Incheon Airport Maglev (opened July 2014)
Changsha Maglev Express (opened 2016)
S1 line of Beijing Subway (opened 2017)
Amusement rides
There are many roller coasters throughout the world that use LIMs to accelerate the ride vehicles. The first being Flight of Fear at Kings Island and Kings Dominion, both opening in 1996. Battlestar Galactica: Human VS Cylon & Revenge of the Mummy at Universal Studios Singapore opened in 2010. They both use LIMs to accelerate from certain point in the rides.
Revenge of the Mummy also located at Universal Studios Hollywood and Universal Studios Florida. The Incredible Hulk Coaster and VelociCoaster at Universal Islands of Adventure also use linear motors. At Walt Disney World, Rock 'n' Roller Coaster Starring Aerosmith at Disney's Hollywood Studios and Guardians of the Galaxy: Cosmic Rewind at Epcot both use LSM to launch their ride vehicles into their indoor ride enclosures.
In 2023 a hydraulic launch roller coaster, Top Thrill Dragster at Cedar Point in Ohio, USA, was renovated and the hydraulic launch replaced with a weaker multi-launch system using LSM, that creates less g-force.
Aircraft launching
Electromagnetic Aircraft Launch System
Proposed and research
Launch loop – A proposed system for launching vehicles into space using a linear motor powered loop
StarTram – Concept for a linear motor on extreme scale
Tether cable catapult system
Aérotrain S44 – A suburban commuter hovertrain prototype
Research Test Vehicle 31 – A hovercraft-type vehicle guided by a track
Hyperloop – a conceptual high-speed transportation system put forward by entrepreneur Elon Musk
Elevator
Lift
Magway - a UK freight delivery system under research and development that aims to deliver goods in pods via 90 cm diameter pipework under and over ground.
| Technology | Engines | null |
44545 | https://en.wikipedia.org/wiki/Till | Till | Till or glacial till is unsorted glacial sediment.
Till is derived from the erosion and entrainment of material by the moving ice of a glacier. It is deposited some distance down-ice to form terminal, lateral, medial and ground moraines.
Till is classified into primary deposits, laid down directly by glaciers, and secondary deposits, reworked by fluvial transport and other processes.
Description
Till is a form of glacial drift, which is rock material transported by a glacier and deposited directly from the ice or from running water emerging from the ice. It is distinguished from other forms of drift in that it is deposited directly by glaciers without being reworked by meltwater.
Till is characteristically unsorted and unstratified, and is not usually consolidated. Most till consists predominantly of clay, silt, and sand, but with pebbles, cobbles, and boulders scattered through the till. The abundance of clay demonstrates lack of reworking by turbulent flow, which otherwise would winnow the clay. Typically, the distribution of particle sizes shows two peaks (it is bimodal) with pebbles predominating in the coarser peak.
The larger clasts (rock fragments) in till typically show a diverse composition, often including rock types from outcrops hundreds of kilometers away. Some clasts may be rounded, and these are thought to be stream pebbles entrained by the glacier. Many of the clasts are faceted, striated, or polished, all signs of glacial abrasion. The sand and silt grains are typically angular to subangular rather than rounded.
It has been known since the careful statistic work by geologist Chauncey D. Holmes in 1941 that elongated clasts in tills tend to align with the direction of ice flow. Clasts in till may also show slight imbrication, with the clasts dipping upstream.
Though till is generally unstratified, till high in clay may show lamination due to compaction under the weight of overlying ice. Till may also contain lenses of sand or gravel, indicating minor and local reworking by water transitional to non-till glacial drift.
The term till comes from an old Scottish name for coarse, rocky soil. It was first used to describe primary glacial deposits by Archibald Geikie in 1863. Early researchers tended to prefer the term boulder clay for the same kind of sediments, but this has fallen into disfavor. Where it is unclear whether a poorly sorted, unconsolidated glacial deposit was deposited directly from glaciers, it is described as diamict or (when lithified) as diamictite. Tillite is a sedimentary rock formed by lithification of till.
Processes
Erosional
Glacial till is mostly derived from subglacial erosion and from the entrainment by the moving ice of previously available unconsolidated sediments. Bedrock can be eroded through the action of glacial plucking and abrasion, and the resulting clasts of various sizes will be incorporated to the glacier's bed.
Glacial abrasion is the weathering of bedrock below a flowing glacier by fragmented rock on the basal layer of the glacier. The two mechanisms of glacial abrasion are striation of the bedrock by coarse grains moved by the glacier, thus gouging the rock below, and polishing of the bedrock by smaller grains such as silts. Glacial plucking is the removal of large blocks from the bed of a glacier.
Much of the silt in till is produced by glacial grinding, and the longer the till remains at the ice-bedrock interface, the more thoroughly it is crushed. However, the crushing process appears to stop with fine silt. Clay in till is likely eroded from bedrock rather than being created by glacial processes.
Depositional
The sediments carried by a glacier will eventually be deposited some distance down-ice from its source. This takes place in the ablation zone, which is the part of the glacier where the rate of ablation (removal of ice by evaporation, melting, or other processes) exceeds the rate of accumulation of new ice from snowfall. As ice is removed, debris are left behind as till. The deposition of glacial till is not uniform, and a single till plain can contain a wide variety of different types of tills due to the various erosional mechanisms and location of till with respect to the transporting glacier.
The different types of till can be categorized between subglacial (beneath) and supraglacial (surface) deposits. Subglacial deposits include lodgement, subglacial meltout, and deformation tills. Supraglacial deposits include supraglacial meltout and flow till. Supraglacial deposits and landforms are widespread in areas of glacial downwasting (vertical thinning of glaciers, as opposed to ice-retreat. They typically sit at the top of the stratigraphic sediment sequence, which has a major influence on land usage. Till is deposited as the terminal moraine, along the lateral and medial moraines and in the ground moraine of a glacier, and moraine is often conflated with till in older writings. Till may also be deposited as drumlins and flutes, though some drumlins consist of a core of stratified sediments with only a cover of till. Interpreting the glacial history of landforms can be difficult due to the tendency of overprinting landforms on top of each other.
As a glacier melts, large amounts of till are eroded and become a source of sediments for reworked glacial drift deposits. These include glaciofluvial deposits, such as outwash in sandurs, and as glaciolacustrine and glaciomarine deposits, such as varves (annual layers) in any proglacial lakes which may form. Erosion of till may take place even in the subglacial environment, such as in tunnel valleys.
Types of till
There are various types of classifying tills:
Primary deposits – Laid down directly by glacier action.
Secondary deposits – Reworked by fluvial transport, erosion, etc.
Traditionally (e.g. Dreimanis, 1988) a further set of divisions has been made to primary deposits, based upon the method of deposition. Van der Meer et al. 2003 have suggested that these till classifications are outdated and should instead be replaced with only one classification, that of deformation till. The reasons behind this are largely down to the difficulties in accurately classifying different tills, which are often based on inferences of the physical setting of the till rather than detailed analysis of the till fabric or particle size.
Subglacial till
Lodgement till
Subglacial lodgement tills are deposits beneath the glacier that are forced, or "lodged" into the bed below. As glaciers advance or retreat, the clasts that are deposited by the ice may have a lower velocity than the ice itself. When the friction between the clast and the bed exceeds the forces of the ice flowing above and around it, the clast will cease to move, and it will become a lodgement till.
Meltout till
Subglacial meltout tills are tills that are deposited via the melting of the ice lobe. Clasts are transported to the base of the glacier over time, and as basal melting continues, they are slowly deposited below the glacier. Since the rate of deposition is controlled by the rate of basal melting, it is worth considering the factors that contribute to melting. These can be the geothermal heat flux, frictional heat generated by sliding, ice thickness, and ice-surface temperature gradients.
Deformation till
Subglacial deformation tills refer to the homogenization of glacial sediments that occur when the stresses and shear forces from the moving glacier rework the topography of the bed. These contain preglacial sediments (non glacial or earlier glacial sediments), which have been run over and thus deformed by meltout processes or lodgement. The constant reworking of these deposited tills leads to a highly homogenized till.
Supraglacial till
Meltout till
Supraglacial meltout tills are similar to subglacial meltout tills. Rather than being the product of basal melting, however, supraglacial meltout tills are imposed on top of the glacier. These consist of clasts and debris that become exposed due to melting via solar radiation. These debris are either just debris that have a high relative position on the glacier, or clasts that have been transported up from the base of the glacier. Debris accumulation has a feedback-loop relationship with melting. Initially, the darker colored debris absorb more heat and thus accelerate the melting process. After a significant amount of melting has occurred, the thickness of the till insulates the ice sheet and slows the melting process. Supraglacial meltout tills typically end up forming moraines.
Flow till
Supraglacial flow tills refer to tills that are subject to a dense concentration of clasts and debris from meltout. These debris localities are then subsequently affected by ablation. Due to their unstable nature, they are subject to downslope flow, and thus named "flow till." Properties of flow tills vary, and can depend on factors such as water content, surface gradient, and debris characteristics. Generally, flow tills with a higher water content behave more fluidly, and thus are more susceptible to flow. There are three main types of flows, which are listed below.
Mobile flows: Thin, fluid, and rapid flows that significantly contribute to erosional processes. These cause strong clast orientation in the direction of flow.
Semi-plastic: Thick, slow moving "tongues" of debris. These are also erosive, and clast sorting is more organized than in mobile flows.
Creep: Very slow movement of debris, downslope in direction. Flow rate is slow enough not to be seen on relatively short timescales, as observed by humans. Particle orientation is often random and not associated with the direction of flow.
Tillite
In cases where till has been indurated or lithified by subsequent burial into solid rock, it is known as the sedimentary rock tillite. Matching beds of ancient tillites on opposite sides of the south Atlantic Ocean provided early evidence for continental drift. The same tillites also provide some support to the Precambrian Snowball Earth glaciation event hypothesis.
Economic resources
Tills sometimes contain placer deposits of valuable minerals such as gold. Diamonds have been found in glacial till in the north-central United States and in Canada. Till prospecting is a method of prospecting in which tills are sampled over a wide area to determine if they contain valuable minerals, such as gold, uranium, silver, nickel, or diamonds, and the flow direction indicated by the till is then used to trace the minerals back to their bedrock source.
| Physical sciences | Sedimentology | Earth science |
44568 | https://en.wikipedia.org/wiki/Herbivore | Herbivore | A herbivore is an animal anatomically and physiologically evolved to feed on plants, especially upon vascular tissues such as foliage, fruits or seeds, as the main component of its diet. These more broadly also encompass animals that eat non-vascular autotrophs such as mosses, algae and lichens, but do not include those feeding on decomposed plant matters (i.e. detritivores) or macrofungi (i.e. fungivores).
As a result of their plant-based diet, herbivorous animals typically have mouth structures (jaws or mouthparts) well adapted to mechanically break down plant materials, and their digestive systems have special enzymes (e.g. amylase and cellulase) to digest polysaccharides. Grazing herbivores such as horses and cattles have wide flat-crowned teeth that are better adapted for grinding grass, tree bark and other tougher lignin-containing materials, and many of them evolved rumination or cecotropic behaviors to better extract nutrients from plants. A large percentage of herbivores also have mutualistic gut flora made up of bacteria and protozoans that help to degrade the cellulose in plants, whose heavily cross-linking polymer structure makes it far more difficult to digest than the protein- and fat-rich animal tissues that carnivores eat.
Etymology
Herbivore is the anglicized form of a modern Latin coinage, herbivora, cited in Charles Lyell's 1830 Principles of Geology. Richard Owen employed the anglicized term in an 1854 work on fossil teeth and skeletons. Herbivora is derived from Latin herba 'small plant, herb' and vora, from vorare 'to eat, devour'.
Definition and related terms
Herbivory is a form of consumption in which an organism principally eats autotrophs such as plants, algae and photosynthesizing bacteria. More generally, organisms that feed on autotrophs in general are known as primary consumers.
Herbivory is usually limited to animals that eat plants. Insect herbivory can cause a variety of physical and metabolic alterations in the way the host plant interacts with itself and other surrounding biotic factors. Fungi, bacteria, and protists that feed on living plants are usually termed plant pathogens (plant diseases), while fungi and microbes that feed on dead plants are described as saprotrophs. Flowering plants that obtain nutrition from other living plants are usually termed parasitic plants. There is, however, no single exclusive and definitive ecological classification of consumption patterns; each textbook has its own variations on the theme.
Evolution of herbivory
The understanding of herbivory in geological time comes from three sources: fossilized plants, which may preserve evidence of defence (such as spines), or herbivory-related damage; the observation of plant debris in fossilised animal faeces; and the construction of herbivore mouthparts.
Although herbivory was long thought to be a Mesozoic phenomenon, fossils have shown that plants were being consumed by arthropods within less than 20 million years after the first land plants evolved. Insects fed on the spores of early Devonian plants, and the Rhynie chert also provides evidence that organisms fed on plants using a "pierce and suck" technique.
During the next 75 million years, plants evolved a range of more complex organs, such as roots and seeds. There is no evidence of any organism being fed upon until the middle-late Mississippian, . There was a gap of 50 to 100 million years between the time each organ evolved and the time organisms evolved to feed upon them; this may be due to the low levels of oxygen during this period, which may have suppressed evolution. Further than their arthropod status, the identity of these early herbivores is uncertain. Hole feeding and skeletonization are recorded in the early Permian, with surface fluid feeding evolving by the end of that period.
Herbivory among four-limbed terrestrial vertebrates, the tetrapods, developed in the Late Carboniferous (307–299 million years ago). The oldest known example being Desmatodon hesperis. Early tetrapods were large amphibious piscivores. While amphibians continued to feed on fish and insects, some reptiles began exploring two new food types, tetrapods (carnivory) and plants (herbivory). The entire dinosaur order ornithischia was composed of herbivorous dinosaurs. Carnivory was a natural transition from insectivory for medium and large tetrapods, requiring minimal adaptation. In contrast, a complex set of adaptations was necessary for feeding on highly fibrous plant materials.
Arthropods evolved herbivory in four phases, changing their approach to it in response to changing plant communities. Tetrapod herbivores made their first appearance in the fossil record of their jaws near the Permio-Carboniferous boundary, approximately 300 million years ago. The earliest evidence of their herbivory has been attributed to dental occlusion, the process in which teeth from the upper jaw come in contact with teeth in the lower jaw is present. The evolution of dental occlusion led to a drastic increase in plant food processing and provides evidence about feeding strategies based on tooth wear patterns. Examination of phylogenetic frameworks of tooth and jaw morphologes has revealed that dental occlusion developed independently in several lineages tetrapod herbivores. This suggests that evolution and spread occurred simultaneously within various lineages.
Food chain
Herbivores form an important link in the food chain because they consume plants to digest the carbohydrates photosynthetically produced by a plant. Carnivores in turn consume herbivores for the same reason, while omnivores can obtain their nutrients from either plants or animals. Due to a herbivore's ability to survive solely on tough and fibrous plant matter, they are termed the primary consumers in the food cycle (chain). Herbivory, carnivory, and omnivory can be regarded as special cases of consumer–resource interactions.
Feeding strategies
Two herbivore feeding strategies are grazing (e.g. cows) and browsing (e.g. moose). For a terrestrial mammal to be called a grazer, at least 90% of the forage has to be grass, and for a browser at least 90% tree leaves and twigs. An intermediate feeding strategy is called "mixed-feeding". In their daily need to take up energy from forage, herbivores of different body mass may be selective in choosing their food. "Selective" means that herbivores may choose their forage source depending on, e.g., season or food availability, but also that they may choose high quality (and consequently highly nutritious) forage before lower quality. The latter especially is determined by the body mass of the herbivore, with small herbivores selecting for high-quality forage, and with increasing body mass animals are less selective. Several theories attempt to explain and quantify the relationship between animals and their food, such as Kleiber's law, Holling's disk equation and the marginal value theorem (see below).
Kleiber's law describes the relationship between an animal's size and its feeding strategy, saying that larger animals need to eat less food per unit weight than smaller animals. Kleiber's law states that the metabolic rate (q0) of an animal is the mass of the animal (M) raised to the 3/4 power: q0=M3/4
Therefore, the mass of the animal increases at a faster rate than the metabolic rate.
Herbivores employ numerous types of feeding strategies. Many herbivores do not fall into one specific feeding strategy, but employ several strategies and eat a variety of plant parts.
Optimal foraging theory is a model for predicting animal behavior while looking for food or other resources, such as shelter or water. This model assesses both individual movement, such as animal behavior while looking for food, and distribution within a habitat, such as dynamics at the population and community level. For example, the model would be used to look at the browsing behavior of a deer while looking for food, as well as that deer's specific location and movement within the forested habitat and its interaction with other deer while in that habitat.
This model has been criticized as circular and untestable. Critics have pointed out that its proponents use examples that fit the theory, but do not use the model when it does not fit the reality. Other critics point out that animals do not have the ability to assess and maximize their potential gains, therefore the optimal foraging theory is irrelevant and derived to explain trends that do not exist in nature.
Holling's disk equation models the efficiency at which predators consume prey. The model predicts that as the number of prey increases, the amount of time predators spend handling prey also increases, and therefore the efficiency of the predator decreases. In 1959, S. Holling proposed an equation to model the rate of return for an optimal diet: Rate (R )=Energy gained in foraging (Ef)/(time searching (Ts) + time handling (Th))
Where s=cost of search per unit time f=rate of encounter with items, h=handling time, e=energy gained per encounter.
In effect, this would indicate that a herbivore in a dense forest would spend more time handling (eating) the vegetation because there was so much vegetation around than a herbivore in a sparse forest, who could easily browse through the forest vegetation. According to the Holling's disk equation, a herbivore in the sparse forest would be more efficient at eating than the herbivore in the dense forest.
The marginal value theorem describes the balance between eating all the food in a patch for immediate energy, or moving to a new patch and leaving the plants in the first patch to regenerate for future use. The theory predicts that absent complicating factors, an animal should leave a resource patch when the rate of payoff (amount of food) falls below the average rate of payoff for the entire area. According to this theory, an animal should move to a new patch of food when the patch they are currently feeding on requires more energy to obtain food than an average patch. Within this theory, two subsequent parameters emerge, the Giving Up Density (GUD) and the Giving Up Time (GUT). The Giving Up Density (GUD) quantifies the amount of food that remains in a patch when a forager moves to a new patch. The Giving Up Time (GUT) is used when an animal continuously assesses the patch quality.
Plant-herbivore interactions
Interactions between plants and herbivores can play a prevalent role in ecosystem dynamics such community structure and functional processes. Plant diversity and distribution is often driven by herbivory, and it is likely that trade-offs between plant competitiveness and defensiveness, and between colonization and mortality allow for coexistence between species in the presence of herbivores. However, the effects of herbivory on plant diversity and richness is variable. For example, increased abundance of herbivores such as deer decrease plant diversity and species richness, while other large mammalian herbivores like bison control dominant species which allows other species to flourish. Plant-herbivore interactions can also operate so that plant communities mediate herbivore communities. Plant communities that are more diverse typically sustain greater herbivore richness by providing a greater and more diverse set of resources.
Coevolution and phylogenetic correlation between herbivores and plants are important aspects of the influence of herbivore and plant interactions on communities and ecosystem functioning, especially in regard to herbivorous insects. This is apparent in the adaptations plants develop to tolerate and/or defend from insect herbivory and the responses of herbivores to overcome these adaptations. The evolution of antagonistic and mutualistic plant-herbivore interactions are not mutually exclusive and may co-occur. Plant phylogeny has been found to facilitate the colonization and community assembly of herbivores, and there is evidence of phylogenetic linkage between plant beta diversity and phylogenetic beta diversity of insect clades such as butterflies. These types of eco-evolutionary feedbacks between plants and herbivores are likely the main driving force behind plant and herbivore diversity.
Abiotic factors such as climate and biogeographical features also impact plant-herbivore communities and interactions. For example, in temperate freshwater wetlands herbivorous waterfowl communities change according to season, with species that eat above-ground vegetation being abundant during summer, and species that forage below-ground being present in winter months. These seasonal herbivore communities differ in both their assemblage and functions within the wetland ecosystem. Such differences in herbivore modalities can potentially lead to trade-offs that influence species traits and may lead to additive effects on community composition and ecosystem functioning. Seasonal changes and environmental gradients such as elevation and latitude often affect the palatability of plants which in turn influences herbivore community assemblages and vice versa. Examples include a decrease in abundance of leaf-chewing larvae in the fall when hardwood leaf palatability decreases due to increased tannin levels which results in a decline of arthropod species richness, and increased palatability of plant communities at higher elevations where grasshoppers abundances are lower. Climatic stressors such as ocean acidification can lead to responses in plant-herbivore interactions in relation to palatability as well.
Herbivore offense
The myriad defenses displayed by plants means that their herbivores need a variety of skills to overcome these defenses and obtain food. These allow herbivores to increase their feeding and use of a host plant. Herbivores have three primary strategies for dealing with plant defenses: choice, herbivore modification, and plant modification.
Feeding choice involves which plants a herbivore chooses to consume. It has been suggested that many herbivores feed on a variety of plants to balance their nutrient uptake and to avoid consuming too much of any one type of defensive chemical. This involves a tradeoff however, between foraging on many plant species to avoid toxins or specializing on one type of plant that can be detoxified.
Herbivore modification is when various adaptations to body or digestive systems of the herbivore allow them to overcome plant defenses. This might include detoxifying secondary metabolites, sequestering toxins unaltered, or avoiding toxins, such as through the production of large amounts of saliva to reduce effectiveness of defenses. Herbivores may also utilize symbionts to evade plant defenses. For example, some aphids use bacteria in their gut to provide essential amino acids lacking in their sap diet.
Plant modification occurs when herbivores manipulate their plant prey to increase feeding. For example, some caterpillars roll leaves to reduce the effectiveness of plant defenses activated by sunlight.
Plant defense
A plant defense is a trait that increases plant fitness when faced with herbivory. This is measured relative to another plant that lacks the defensive trait. Plant defenses increase survival and/or reproduction (fitness) of plants under pressure of predation from herbivores.
Defense can be divided into two main categories, tolerance and resistance. Tolerance is the ability of a plant to withstand damage without a reduction in fitness. This can occur by diverting herbivory to non-essential plant parts, resource allocation, compensatory growth, or by rapid regrowth and recovery from herbivory. Resistance refers to the ability of a plant to reduce the amount of damage it receives from herbivores. This can occur via avoidance in space or time, physical defenses, or chemical defenses. Defenses can either be constitutive, always present in the plant, or induced, produced or translocated by the plant following damage or stress.
Physical, or mechanical, defenses are barriers or structures designed to deter herbivores or reduce intake rates, lowering overall herbivory. Thorns such as those found on roses or acacia trees are one example, as are the spines on a cactus. Smaller hairs known as trichomes may cover leaves or stems and are especially effective against invertebrate herbivores. In addition, some plants have waxes or resins that alter their texture, making them difficult to eat. Also the incorporation of silica into cell walls is analogous to that of the role of lignin in that it is a compression-resistant structural component of cell walls; so that plants with their cell walls impregnated with silica are thereby afforded a measure of protection against herbivory.
Chemical defenses are secondary metabolites produced by the plant that deter herbivory. There are a wide variety of these in nature and a single plant can have hundreds of different chemical defenses. Chemical defenses can be divided into two main groups, carbon-based defenses and nitrogen-based defenses.
Carbon-based defenses include terpenes and phenolics. Terpenes are derived from 5-carbon isoprene units and comprise essential oils, carotenoids, resins, and latex. They can have several functions that disrupt herbivores such as inhibiting adenosine triphosphate (ATP) formation, molting hormones, or the nervous system. Phenolics combine an aromatic carbon ring with a hydroxyl group. There are several different phenolics such as lignins, which are found in cell walls and are very indigestible except for specialized microorganisms; tannins, which have a bitter taste and bind to proteins making them indigestible; and furanocumerins, which produce free radicals disrupting DNA, protein, and lipids, and can cause skin irritation.
Nitrogen-based defenses are synthesized from amino acids and primarily come in the form of alkaloids and cyanogens. Alkaloids include commonly recognized substances such as caffeine, nicotine, and morphine. These compounds are often bitter and can inhibit DNA or RNA synthesis or block nervous system signal transmission. Cyanogens get their name from the cyanide stored within their tissues. This is released when the plant is damaged and inhibits cellular respiration and electron transport.
Plants have also changed features that enhance the probability of attracting natural enemies to herbivores. Some emit semiochemicals, odors that attract natural enemies, while others provide food and housing to maintain the natural enemies' presence, e.g. ants that reduce herbivory. A given plant species often has many types of defensive mechanisms, mechanical or chemical, constitutive or induced, which allow it to escape from herbivores.
Predator–prey theory
According to the theory of predator–prey interactions, the relationship between herbivores and plants is cyclic. When prey (plants) are numerous their predators (herbivores) increase in numbers, reducing the prey population, which in turn causes predator number to decline. The prey population eventually recovers, starting a new cycle. This suggests that the population of the herbivore fluctuates around the carrying capacity of the food source, in this case, the plant.
Several factors play into these fluctuating populations and help stabilize predator-prey dynamics. For example, spatial heterogeneity is maintained, which means there will always be pockets of plants not found by herbivores. This stabilizing dynamic plays an especially important role for specialist herbivores that feed on one species of plant and prevents these specialists from wiping out their food source. Prey defenses also help stabilize predator-prey dynamics, and for more information on these relationships see the section on Plant Defenses. Eating a second prey type helps herbivores' populations stabilize. Alternating between two or more plant types provides population stability for the herbivore, while the populations of the plants oscillate. This plays an important role for generalist herbivores that eat a variety of plants. Keystone herbivores keep vegetation populations in check and allow for a greater diversity of both herbivores and plants. When an invasive herbivore or plant enters the system, the balance is thrown off and the diversity can collapse to a monotaxon system.
The back and forth relationship of plant defense and herbivore offense drives coevolution between plants and herbivores, resulting in a "coevolutionary arms race". The escape and radiation mechanisms for coevolution, presents the idea that adaptations in herbivores and their host plants, has been the driving force behind speciation.
Mutualism
While much of the interaction of herbivory and plant defense is negative, with one individual reducing the fitness of the other, some is beneficial. This beneficial herbivory takes the form of mutualisms in which both partners benefit in some way from the interaction. Seed dispersal by herbivores and pollination are two forms of mutualistic herbivory in which the herbivore receives a food resource and the plant is aided in reproduction. Plants can also be indirectly affected by herbivores through nutrient recycling, with plants benefiting from herbivores when nutrients are recycled very efficiently. Another form of plant-herbivore mutualism is physical changes to the environment and/or plant community structure by herbivores which serve as ecosystem engineers, such as wallowing by bison. Swans form a mutual relationship with the plant species that they forage by digging and disturbing the sediment which removes competing plants and subsequently allows colonization of other plant species.
Impacts
Trophic cascades and environmental degradation
When herbivores are affected by trophic cascades, plant communities can be indirectly affected. Often these effects are felt when predator populations decline and herbivore populations are no longer limited, which leads to intense herbivore foraging which can suppress plant communities. With the size of herbivores having an effect on the amount of energy intake that is needed, larger herbivores need to forage on higher quality or more plants to gain the optimal amount of nutrients and energy compared to smaller herbivores. Environmental degradation from white-tailed deer (Odocoileus virginianus) in the US alone has the potential to both change vegetative communities through over-browsing and cost forest restoration projects upwards of $750 million annually. Another example of a trophic cascade involved plant-herbivore interactions are coral reef ecosystems. Herbivorous fish and marine animals are important algae and seaweed grazers, and in the absence of plant-eating fish, corals are outcompeted and seaweeds deprive corals of sunlight.
Economic impacts
Agricultural crop damage by the same species totals approximately $100 million every year. Insect crop damages also contribute largely to annual crop losses in the U.S. Herbivores also affect economics through the revenue generated by hunting and ecotourism. For example, the hunting of herbivorous game species such as white-tailed deer, cottontail rabbits, antelope, and elk in the U.S. contributes greatly to the billion-dollar annually, hunting industry. Ecotourism is a major source of revenue, particularly in Africa, where many large mammalian herbivores such as elephants, zebras, and giraffes help to bring in the equivalent of millions of US dollars to various nations annually.
| Biology and health sciences | Ethology | null |
44578 | https://en.wikipedia.org/wiki/Big%20O%20notation | Big O notation | Big O notation is a mathematical notation that describes the limiting behavior of a function when the argument tends towards a particular value or infinity. Big O is a member of a family of notations invented by German mathematicians Paul Bachmann, Edmund Landau, and others, collectively called Bachmann–Landau notation or asymptotic notation. The letter O was chosen by Bachmann to stand for Ordnung, meaning the order of approximation.
In computer science, big O notation is used to classify algorithms according to how their run time or space requirements grow as the input size grows. In analytic number theory, big O notation is often used to express a bound on the difference between an arithmetical function and a better understood approximation; a famous example of such a difference is the remainder term in the prime number theorem. Big O notation is also used in many other fields to provide similar estimates.
Big O notation characterizes functions according to their growth rates: different functions with the same asymptotic growth rate may be represented using the same O notation. The letter O is used because the growth rate of a function is also referred to as the order of the function. A description of a function in terms of big O notation usually only provides an upper bound on the growth rate of the function.
Associated with big O notation are several related notations, using the symbols , and , to describe other kinds of bounds on asymptotic growth rates.
Formal definition
Let the function to be estimated, be a real or complex valued function, and let the comparison function, be a real valued function. Let both functions be defined on some unbounded subset of the positive real numbers, and be non-zero (often, but not necessarily, strictly positive) for all large enough values of One writes
and it is read " is big O of " or more often " is of the order of " if the absolute value of is at most a positive constant multiple of the absolute value of for all sufficiently large values of That is, if there exists a positive real number and a real number such that
In many contexts, the assumption that we are interested in the growth rate as the variable goes to infinity or to zero is left unstated, and one writes more simply that
The notation can also be used to describe the behavior of near some real number (often, ): we say
if there exist positive numbers and such that for all defined with
As is non-zero for adequately large (or small) values of both of these definitions can be unified using the limit superior:
if
And in both of these definitions the limit point (whether or not) is a cluster point of the domains of and i. e., in every neighbourhood of there have to be infinitely many points in common. Moreover, as pointed out in the article about the limit inferior and limit superior, the (at least on the extended real number line) always exists.
In computer science, a slightly more restrictive definition is common: and are both required to be functions from some unbounded subset of the positive integers to the nonnegative real numbers; then if there exist positive integer numbers and such that for all
Example
In typical usage the notation is asymptotical, that is, it refers to very large . In this setting, the contribution of the terms that grow "most quickly" will eventually make the other ones irrelevant. As a result, the following simplification rules can be applied:
If is a sum of several terms, if there is one with largest growth rate, it can be kept, and all others omitted.
If is a product of several factors, any constants (factors in the product that do not depend on ) can be omitted.
For example, let , and suppose we wish to simplify this function, using notation, to describe its growth rate as approaches infinity. This function is the sum of three terms: , , and . Of these three terms, the one with the highest growth rate is the one with the largest exponent as a function of , namely . Now one may apply the second rule: is a product of and in which the first factor does not depend on . Omitting this factor results in the simplified form . Thus, we say that is a "big O" of . Mathematically, we can write . One may confirm this calculation using the formal definition: let and . Applying the formal definition from above, the statement that is equivalent to its expansion,
for some suitable choice of a real number and a positive real number and for all . To prove this, let and . Then, for all :
so
Use
Big O notation has two main areas of application:
In mathematics, it is commonly used to describe how closely a finite series approximates a given function, especially in the case of a truncated Taylor series or asymptotic expansion.
In computer science, it is useful in the analysis of algorithms.
In both applications, the function appearing within the is typically chosen to be as simple as possible, omitting constant factors and lower order terms.
There are two formally close, but noticeably different, usages of this notation:
infinite asymptotics
infinitesimal asymptotics.
This distinction is only in application and not in principle, however—the formal definition for the "big O" is the same for both cases, only with different limits for the function argument.
Infinite asymptotics
Big O notation is useful when analyzing algorithms for efficiency. For example, the time (or the number of steps) it takes to complete a problem of size might be found to be . As grows large, the term will come to dominate, so that all other terms can be neglected—for instance when , the term is 1000 times as large as the term. Ignoring the latter would have negligible effect on the expression's value for most purposes. Further, the coefficients become irrelevant if we compare to any other order of expression, such as an expression containing a term or . Even if , if , the latter will always exceed the former once grows larger than , viz. . Additionally, the number of steps depends on the details of the machine model on which the algorithm runs, but different types of machines typically vary by only a constant factor in the number of steps needed to execute an algorithm. So the big O notation captures what remains: we write either
or
and say that the algorithm has order of time complexity. The sign "" is not meant to express "is equal to" in its normal mathematical sense, but rather a more colloquial "is", so the second expression is sometimes considered more accurate (see the "Equals sign" discussion below) while the first is considered by some as an abuse of notation.
Infinitesimal asymptotics
Big O can also be used to describe the error term in an approximation to a mathematical function. The most significant terms are written explicitly, and then the least-significant terms are summarized in a single big O term. Consider, for example, the exponential series and two expressions of it that are valid when is small:
The middle expression (the one with O(x3)) means the absolute-value of the error ex − (1 + x + x2/2) is at most some constant times x3 when x is close enough to 0.
Properties
If the function can be written as a finite sum of other functions, then the fastest growing one determines the order of . For example,
In particular, if a function may be bounded by a polynomial in , then as tends to infinity, one may disregard lower-order terms of the polynomial. The sets and are very different. If is greater than one, then the latter grows much faster. A function that grows faster than for any is called superpolynomial. One that grows more slowly than any exponential function of the form is called subexponential. An algorithm can require time that is both superpolynomial and subexponential; examples of this include the fastest known algorithms for integer factorization and the function .
We may ignore any powers of inside of the logarithms. The set is exactly the same as . The logarithms differ only by a constant factor (since ) and thus the big O notation ignores that. Similarly, logs with different constant bases are equivalent. On the other hand, exponentials with different bases are not of the same order. For example, and are not of the same order.
Changing units may or may not affect the order of the resulting algorithm. Changing units is equivalent to multiplying the appropriate variable by a constant wherever it appears. For example, if an algorithm runs in the order of , replacing by means the algorithm runs in the order of , and the big O notation ignores the constant . This can be written as . If, however, an algorithm runs in the order of , replacing with gives . This is not equivalent to in general. Changing variables may also affect the order of the resulting algorithm. For example, if an algorithm's run time is when measured in terms of the number of digits of an input number , then its run time is when measured as a function of the input number itself, because .
Product
Sum
If and then . It follows that if and then . In other words, this second statement says that is a convex cone.
Multiplication by a constant
Let be a nonzero constant. Then . In other words, if , then
Multiple variables
Big O (and little o, Ω, etc.) can also be used with multiple variables. To define big O formally for multiple variables, suppose and are two functions defined on some subset of . We say
if and only if there exist constants and such that for all with for some
Equivalently, the condition that for some can be written , where denotes the Chebyshev norm. For example, the statement
asserts that there exist constants C and M such that
whenever either or holds. This definition allows all of the coordinates of to increase to infinity. In particular, the statement
(i.e., ) is quite different from
(i.e., ).
Under this definition, the subset on which a function is defined is significant when generalizing statements from the univariate setting to the multivariate setting. For example, if and , then if we restrict and to , but not if they are defined on .
This is not the only generalization of big O to multivariate functions, and in practice, there is some inconsistency in the choice of definition.
Matters of notation
Equals sign
The statement " is " as defined above is usually written as . Some consider this to be an abuse of notation, since the use of the equals sign could be misleading as it suggests a symmetry that this statement does not have. As de Bruijn says, is true but is not. Knuth describes such statements as "one-way equalities", since if the sides could be reversed, "we could deduce ridiculous things like from the identities and ". In another letter, Knuth also pointed out that
"the equality sign is not symmetric with respect to such notations", [as, in this notation,] "mathematicians customarily use the '=' sign as they use the word 'is' in English: Aristotle is a man, but a man isn't necessarily Aristotle".
For these reasons, it would be more precise to use set notation and write (read as: " is an element of ", or " is in the set thinking of as the class of all functions such that for some positive real number . However, the use of the equals sign is customary.
Other arithmetic operators
Big O notation can also be used in conjunction with other arithmetic operators in more complicated equations. For example, denotes the collection of functions having the growth of h(x) plus a part whose growth is limited to that of f(x). Thus,
expresses the same as
Example
Suppose an algorithm is being developed to operate on a set of n elements. Its developers are interested in finding a function T(n) that will express how long the algorithm will take to run (in some arbitrary measurement of time) in terms of the number of elements in the input set. The algorithm works by first calling a subroutine to sort the elements in the set and then perform its own operations. The sort has a known time complexity of O(n2), and after the subroutine runs the algorithm must take an additional steps before it terminates. Thus the overall time complexity of the algorithm can be expressed as . Here the terms are subsumed within the faster-growing O(n2). Again, this usage disregards some of the formal meaning of the "=" symbol, but it does allow one to use the big O notation as a kind of convenient placeholder.
Multiple uses
In more complicated usage, O(·) can appear in different places in an equation, even several times on each side. For example, the following are true for :
The meaning of such statements is as follows: for any functions which satisfy each O(·) on the left side, there are some functions satisfying each O(·) on the right side, such that substituting all these functions into the equation makes the two sides equal. For example, the third equation above means: "For any function f(n) = O(1), there is some function g(n) = O(en) such that nf(n) = g(n)." In terms of the "set notation" above, the meaning is that the class of functions represented by the left side is a subset of the class of functions represented by the right side. In this use the "=" is a formal symbol that unlike the usual use of "=" is not a symmetric relation. Thus for example does not imply the false statement .
Typesetting
Big O is typeset as an italicized uppercase "O", as in the following example: . In TeX, it is produced by simply typing O inside math mode. Unlike Greek-named Bachmann–Landau notations, it needs no special symbol. However, some authors use the calligraphic variant instead.
Orders of common functions
Here is a list of classes of functions that are commonly encountered when analyzing the running time of an algorithm. In each case, c is a positive constant and n increases without bound. The slower-growing functions are generally listed first.
The statement is sometimes weakened to to derive simpler formulas for asymptotic complexity. For any and is a subset of for any so may be considered as a polynomial with some bigger order.
Related asymptotic notations
Big O is widely used in computer science. Together with some other related notations, it forms the family of Bachmann–Landau notations.
Little-o notation
Intuitively, the assertion " is " (read " is little-o of " or " is of inferior order to ") means that grows much faster than , or equivalently grows much slower than . As before, let f be a real or complex valued function and g a real valued function, both defined on some unbounded subset of the positive real numbers, such that g(x) is strictly positive for all large enough values of x. One writes
if for every positive constant there exists a constant such that
For example, one has
and both as
The difference between the definition of the big-O notation and the definition of little-o is that while the former has to be true for at least one constant M, the latter must hold for every positive constant , however small. In this way, little-o notation makes a stronger statement than the corresponding big-O notation: every function that is little-o of g is also big-O of g, but not every function that is big-O of g is little-o of g. For example, but
If g(x) is nonzero, or at least becomes nonzero beyond a certain point, the relation is equivalent to
(and this is in fact how Landau originally defined the little-o notation).
Little-o respects a number of arithmetic operations. For example,
if is a nonzero constant and then , and
if and then
It also satisfies a transitivity relation:
if and then
Big Omega notation
Another asymptotic notation is , read "big omega". There are two widespread and incompatible definitions of the statement
as ,
where a is some real number, , or , where f and g are real functions defined in a neighbourhood of a, and where g is positive in this neighbourhood.
The Hardy–Littlewood definition is used mainly in analytic number theory, and the Knuth definition mainly in computational complexity theory; the definitions are not equivalent.
The Hardy–Littlewood definition
In 1914 G.H. Hardy and J.E. Littlewood introduced the new symbol which is defined as follows:
as if
Thus is the negation of
In 1916 the same authors introduced the two new symbols and defined as:
as if
as if
These symbols were used by E. Landau, with the same meanings, in 1924. Authors that followed Landau, however, use a different notation for the same definitions: The symbol has been replaced by the current notation with the same definition, and became
These three symbols as well as (meaning that and are both satisfied), are now currently used in analytic number theory.
Simple examples
We have
as
and more precisely
as
We have
as
and more precisely
as
however
as
The Knuth definition
In 1976 Donald Knuth published a paper to justify his use of the -symbol to describe a stronger property. Knuth wrote: "For all the applications I have seen so far in computer science, a stronger requirement ... is much more appropriate". He defined
with the comment: "Although I have changed Hardy and Littlewood's definition of , I feel justified in doing so because their definition is by no means in wide use, and because there are other ways to say what they want to say in the comparatively rare cases when their definition applies."
Family of Bachmann–Landau notations
The limit definitions assume for sufficiently large . The table is (partly) sorted from smallest to largest, in the sense that (Knuth's version of) on functions correspond to on the real line (the Hardy–Littlewood version of , however, doesn't correspond to any such description).
Computer science uses the big , big Theta , little , little omega and Knuth's big Omega notations. Analytic number theory often uses the big , small , Hardy's , Hardy–Littlewood's big Omega (with or without the +, − or ± subscripts) and notations. The small omega notation is not used as often in analysis.
Use in computer science
Informally, especially in computer science, the big O notation often can be used somewhat differently to describe an asymptotic tight bound where using big Theta Θ notation might be more factually appropriate in a given context. For example, when considering a function T(n) = 73n3 + 22n2 + 58, all of the following are generally acceptable, but tighter bounds (such as numbers 2 and 3 below) are usually strongly preferred over looser bounds (such as number 1 below).
The equivalent English statements are respectively:
T(n) grows asymptotically no faster than n100
T(n) grows asymptotically no faster than n3
T(n) grows asymptotically as fast as n3.
So while all three statements are true, progressively more information is contained in each. In some fields, however, the big O notation (number 2 in the lists above) would be used more commonly than the big Theta notation (items numbered 3 in the lists above). For example, if T(n) represents the running time of a newly developed algorithm for input size n, the inventors and users of the algorithm might be more inclined to put an upper asymptotic bound on how long it will take to run without making an explicit statement about the lower asymptotic bound.
Other notation
In their book Introduction to Algorithms, Cormen, Leiserson, Rivest and Stein consider the set of functions f which satisfy
In a correct notation this set can, for instance, be called O(g), where
The authors state that the use of equality operator (=) to denote set membership rather than the set membership operator (∈) is an abuse of notation, but that doing so has advantages. Inside an equation or inequality, the use of asymptotic notation stands for an anonymous function in the set O(g), which eliminates lower-order terms, and helps to reduce inessential clutter in equations, for example:
Extensions to the Bachmann–Landau notations
Another notation sometimes used in computer science is Õ (read soft-O), which hides polylogarithmic factors. There are two definitions in use: some authors use f(n) = Õ(g(n)) as shorthand for for some k, while others use it as shorthand for . When is polynomial in n, there is no difference; however, the latter definition allows one to say, e.g. that while the former definition allows for for any constant k. Some authors write O* for the same purpose as the latter definition. Essentially, it is big O notation, ignoring logarithmic factors because the growth-rate effects of some other super-logarithmic function indicate a growth-rate explosion for large-sized input parameters that is more important to predicting bad run-time performance than the finer-point effects contributed by the logarithmic-growth factor(s). This notation is often used to obviate the "nitpicking" within growth-rates that are stated as too tightly bounded for the matters at hand (since logk n is always o(nε) for any constant k and any ).
Also, the L notation, defined as
is convenient for functions that are between polynomial and exponential in terms of
Generalizations and related usages
The generalization to functions taking values in any normed vector space is straightforward (replacing absolute values by norms), where f and g need not take their values in the same space. A generalization to functions g taking values in any topological group is also possible.
The "limiting process" can also be generalized by introducing an arbitrary filter base, i.e. to directed nets f and g. The o notation can be used to define derivatives and differentiability in quite general spaces, and also (asymptotical) equivalence of functions,
which is an equivalence relation and a more restrictive notion than the relationship "f is Θ(g)" from above. (It reduces to lim f / g = 1 if f and g are positive real valued functions.) For example, 2x is Θ(x), but is not o(x).
History (Bachmann–Landau, Hardy, and Vinogradov notations)
The symbol O was first introduced by number theorist Paul Bachmann in 1894, in the second volume of his book Analytische Zahlentheorie ("analytic number theory"). The number theorist Edmund Landau adopted it, and was thus inspired to introduce in 1909 the notation o; hence both are now called Landau symbols. These notations were used in applied mathematics during the 1950s for asymptotic analysis.
The symbol (in the sense "is not an o of") was introduced in 1914 by Hardy and Littlewood. Hardy and Littlewood also introduced in 1916 the symbols ("right") and ("left"), precursors of the modern symbols ("is not smaller than a small o of") and ("is not larger than a small o of"). Thus the Omega symbols (with their original meanings) are sometimes also referred to as "Landau symbols". This notation became commonly used in number theory at least since the 1950s.
The symbol , although it had been used before with different meanings, was given its modern definition by Landau in 1909 and by Hardy in 1910. Just above on the same page of his tract Hardy defined the symbol , where means that both and are satisfied. The notation is still currently used in analytic number theory. In his tract Hardy also proposed the symbol , where means that for some constant .
In the 1970s the big O was popularized in computer science by Donald Knuth, who proposed the different notation for Hardy's , and proposed a different definition for the Hardy and Littlewood Omega notation.
Two other symbols coined by Hardy were (in terms of the modern O notation)
and
(Hardy however never defined or used the notation , nor , as it has been sometimes reported).
Hardy introduced the symbols and (as well as the already mentioned other symbols) in his 1910 tract "Orders of Infinity", and made use of them only in three papers (1910–1913). In his nearly 400 remaining papers and books he consistently used the Landau symbols O and o.
Hardy's symbols and (as well as ) are not used anymore. On the other hand, in the 1930s, the Russian number theorist Ivan Matveyevich Vinogradov introduced his notation , which has been increasingly used in number theory instead of the notation. We have
and frequently both notations are used in the same paper.
The big-O originally stands for "order of" ("Ordnung", Bachmann 1894), and is thus a Latin letter. Neither Bachmann nor Landau ever call it "Omicron". The symbol was much later on (1976) viewed by Knuth as a capital omicron, probably in reference to his definition of the symbol Omega. The digit zero should not be used.
| Mathematics | Algorithms | null |
44593 | https://en.wikipedia.org/wiki/Rib | Rib | In vertebrate anatomy, ribs () are the long curved bones which form the rib cage, part of the axial skeleton. In most tetrapods, ribs surround the thoracic cavity, enabling the lungs to expand and thus facilitate breathing by expanding the thoracic cavity. They serve to protect the lungs, heart, and other vital organs of the thorax. In some animals, especially snakes, ribs may provide support and protection for the entire body.
Human anatomy
Rib details
Human ribs are flat bones that form part of the rib cage to help protect internal organs. Humans usually have 24 ribs, in 12 pairs. 1 in 500 people have an extra rib known as a cervical rib. People may have a cervical rib on the right, left or both sides. All are attached at the back to the thoracic vertebrae and are numbered from 1 to 12 according to the vertebrae to which they attach. The first rib is attached to thoracic vertebra 1 (T1). At the front of the body, most of the ribs are joined by costal cartilage to the sternum. Ribs connect to vertebrae at the costovertebral joints.
The parts of a rib includes the head, neck, body (or shaft), tubercle, and angle.
The head of the rib lies next to a vertebra. The ribs connect to the vertebrae with two costovertebral joints, one on the head and one on the neck. The head of the rib has a superior and an inferior articulating region, separated by a crest. These articulate with the superior and inferior costal facets on the connecting vertebrae. The crest gives attachment to the intra-articulate ligament that joins the rib to the vertebra of the same number, at the intervertebral disc. Another ligament, the radiate ligament joins the head of the rib to both the body of the upper vertebra and to the body of the lower vertebra. The smaller middle part of the ligament connects to the intervertebral disc. This plane joint is known as the articulation of the head of the rib.
The other costovertebral joint is that between the tubercle on the neck and the transverse process of the corresponding thoracic vertebra, known as the costotransverse joint. The superior costotransverse ligament attaches from the non-articular facet of the tubercle to the transverse process of the vertebra.
The neck of the rib is a flattened part that extends laterally from the head. The neck is about 3 cm long. Its anterior surface is flat and smooth, whilst its posterior is perforated by numerous foramina and its surface rough, to give attachment to the ligament of the neck. Its upper border presents a rough crest (crista colli costae) for the attachment of the anterior costotransverse ligament; its lower border is rounded.
A tubercle of rib on the posterior surface of the neck of the rib, has two facets (surfaces) one articulating and one non-articulating. The articular facet, is small and oval and is the lower and more medial of the two, and connects to the transverse costal facet on the thoracic vertebra of the same rib number. The transverse costal facet is on the end of the transverse process of the lower of the two vertebrae to which the head is connected. The non-articular portion is a rough elevation and affords attachment to the ligament of the tubercle. The tubercle is much more prominent in the upper ribs than in the lower ribs.
Rib cage
The first seven sets of ribs, known as "true ribs", are attached to the sternum by the costal cartilages. The first rib is unique and easier to distinguish than other ribs. It is a short, flat, C-shaped bone, and attaches to the manubrium. The vertebral attachment can be found just below the neck at the first thoracic vertebra, and the majority of this bone can be found above the level of the clavicle. Ribs 2 through 7 then become longer and less curved as they progress downwards. The following five sets are known as "false ribs", three of these sharing a common cartilaginous connection to the sternum, while the last two (eleventh and twelfth ribs) are termed floating ribs. They are attached to the vertebrae only, and not to the sternum or cartilage coming off of the sternum.
In general, human ribs increase in length from ribs 1 through 7 and decrease in length again through rib 12. Along with this change in size, the ribs become progressively oblique (slanted) from ribs 1 through 9, then less slanted through rib 12.
The rib cage is separated from the lower abdomen by the thoracic diaphragm which controls breathing. When the diaphragm contracts, the thoracic cavity is expanded, reducing intra-thoracic pressure and drawing air into the lungs. This happens through one of two actions (or a mix of the two): when the lower ribs the diaphragm connects to are stabilized by muscles and the central tendon is mobile, when the muscle contracts the central tendon is drawn down, compressing the cavity underneath and expanding the thoracic cavity downward. When the central tendon is stabilized and the lower ribs are mobile, a contraction of the diaphragm elevates the ribs, which works in conjunction with other muscles to expand the thoracic indent upward.
Development
Early in the developing embryo, somites form and soon subdivide into three mesodermal components – the myotome, dermatome, and the sclerotome. The vertebrae and ribs develop from the sclerotomes.
During the fourth week (fertilization age) costal processes have formed on the vertebral bodies. These processes are small, lateral protrusions of mesenchyme that develop in association with the vertebral arches. During the fifth week the costal processes on the thoracic vertebrae become longer to form the ribs. In the sixth week, the costovertebral joints begin to develop and separate the ribs from the vertebrae. The first seven pairs of ribs, the true ribs join at the front to the sternal bars. By the fetal stage the sternal bars have completely fused.
The ribs begin as cartilage that later ossifies – a process called endochondral ossification. Primary ossification centers are located near the angle of each rib, and ossification continues in the direction away from the head and neck. During adolescence secondary ossification centers are formed in the tubercles and heads of the ribs.
Other animals
In jawed fish, there are often two sets of ribs attached to the vertebral column. One set, the dorsal ribs, are found in the dividing septum between the upper and lower parts of the main muscle segments, projecting roughly sideways from the vertebral column. The second set, the ventral ribs arise from the vertebral column just below the dorsal ribs, and enclose the lower body, often joining at the tips. Not all species possess both types of rib, with the dorsal ribs being most commonly absent. Sharks, for example, have no ventral ribs, and only very short dorsal ribs. In some teleosts, there may be additional rib-like bones within the muscle mass.
Tetrapods, however, only ever have a single set of ribs which are probably homologous with the dorsal ribs of fishes. In the earlier choanates, every vertebra bore a pair of ribs, although those on the thoracic vertebrae are typically the longest. The sacral ribs were stout and short, since they formed part of the pelvis, connecting the backbone to the hip bones.
In most true tetrapods, many of these early ribs have been lost, and in living amphibians and reptiles, there is great variation in rib structure and number. For example, turtles have only eight pairs of ribs, which are developed into a bony or cartilaginous carapace and plastron, while snakes have numerous ribs running along the full length of their trunk. Frogs typically have no ribs, aside from a sacral pair, which form part of the pelvis.
In birds, ribs are present as distinct bones only on the thoracic region, although small fused ribs are present on the cervical vertebrae. The thoracic ribs of birds possess a wide projection to the rear; this uncinate process is an attachment for the shoulder muscles.
Usually dogs have 26 ribs.
Mammals usually also only have distinct ribs on the thoracic vertebra, although fixed cervical ribs are also present in monotremes. In therian mammals, the cervical and lumbar ribs are found only as tiny remnants fused to the vertebrae, where they are referred to as transverse processes. In general, the structure and number of the true ribs in humans is similar to that in other mammals. Unlike reptiles, caudal ribs are never found in mammals.
Ribs as food
Ribs as food are widely used from many animals. The ribs are the less meaty part of the meat chop and they are often cooked as part of a slab; five or more is known as a rack, as in a rack of lamb. Short ribs are ribs of beef either served singly or several as a plate. A rib steak from beef is a popular choice used in many cuisines. Pork ribs, including spare ribs are popular in European and Asian cuisine.
Animated images
| Biology and health sciences | Skeletal system | Biology |
44596 | https://en.wikipedia.org/wiki/Positronium | Positronium | Positronium (Ps) is a system consisting of an electron and its anti-particle, a positron, bound together into an exotic atom, specifically an onium. Unlike hydrogen, the system has no protons. The system is unstable: the two particles annihilate each other to predominantly produce two or three gamma-rays, depending on the relative spin states. The energy levels of the two particles are similar to that of the hydrogen atom (which is a bound state of a proton and an electron). However, because of the reduced mass, the frequencies of the spectral lines are less than half of those for the corresponding hydrogen lines.
States
The mass of positronium is 1.022 MeV, which is twice the electron mass minus the binding energy of a few eV. The lowest energy orbital state of positronium is 1S, and like with hydrogen, it has a hyperfine structure arising from the relative orientations of the spins of the electron and the positron.
The singlet state, , with antiparallel spins (S = 0, Ms = 0) is known as para-positronium (p-Ps). It has a mean lifetime of and decays preferentially into two gamma rays with energy of each (in the center-of-mass frame). Para-positronium can decay into any even number of photons (2, 4, 6, ...), but the probability quickly decreases with the number: the branching ratio for decay into 4 photons is .
Para-positronium lifetime in vacuum is approximately
The triplet states, 3S1, with parallel spins (S = 1, Ms = −1, 0, 1) are known as ortho-positronium (o-Ps), and have an energy that is approximately 0.001 eV higher than the singlet. These states have a mean lifetime of , and the leading decay is three gammas. Other modes of decay are negligible; for instance, the five-photons mode has branching ratio of ≈.
Ortho-positronium lifetime in vacuum can be calculated approximately as:
However more accurate calculations with corrections to O(α2) yield a value of −1 for the decay rate, corresponding to a lifetime of .
Positronium in the 2S state is metastable having a lifetime of against annihilation. The positronium created in such an excited state will quickly cascade down to the ground state, where annihilation will occur more quickly.
Measurements
Measurements of these lifetimes and energy levels have been used in precision tests of quantum electrodynamics, confirming quantum electrodynamics (QED) predictions to high precision.
Annihilation can proceed via a number of channels, each producing gamma rays with total energy of (sum of the electron and positron mass-energy), usually 2 or 3, with up to 5 gamma ray photons recorded from a single annihilation.
The annihilation into a neutrino–antineutrino pair is also possible, but the probability is predicted to be negligible. The branching ratio for o-Ps decay for this channel is (electron neutrino–antineutrino pair) and (for other flavour) in predictions based on the Standard Model, but it can be increased by non-standard neutrino properties, like relatively high magnetic moment. The experimental upper limits on branching ratio for this decay (as well as for a decay into any "invisible" particles) are < for p-Ps and < for o-Ps.
Energy levels
While precise calculation of positronium energy levels uses the Bethe–Salpeter equation or the Breit equation, the similarity between positronium and hydrogen allows a rough estimate. In this approximation, the energy levels are different because of a different effective mass, μ, in the energy equation (see electron energy levels for a derivation):
where:
is the charge magnitude of the electron (same as the positron),
is the Planck constant,
is the electric constant (otherwise known as the permittivity of free space),
is the reduced mass: where and are, respectively, the mass of the electron and the positron (which are the same by definition as antiparticles).
Thus, for positronium, its reduced mass only differs from the electron by a factor of 2. This causes the energy levels to also roughly be half of what they are for the hydrogen atom.
So finally, the energy levels of positronium are given by
The lowest energy level of positronium () is . The next level is . The negative sign is a convention that implies a bound state. Positronium can also be considered by a particular form of the two-body Dirac equation; Two particles with a Coulomb interaction can be exactly separated in the (relativistic) center-of-momentum frame and the resulting ground-state energy has been obtained very accurately using finite element methods of Janine Shertzer. Their results lead to the discovery of anomalous states.
The Dirac equation whose Hamiltonian comprises two Dirac particles and a static Coulomb potential is not relativistically invariant. But if one adds the (or , where is the fine-structure constant) terms, where , then the result is relativistically invariant. Only the leading term is included. The contribution is the Breit term; workers rarely go to because at one has the Lamb shift, which requires quantum electrodynamics.
Formation and decay in materials
After a radioactive atom in a material undergoes a β+ decay (positron emission), the resulting high-energy positron slows down by colliding with atoms, and eventually annihilates with one of the many electrons in the material. It may however first form positronium before the annihilation event. The understanding of this process is of some importance in positron emission tomography. Approximately:
~60% of positrons will directly annihilate with an electron without forming positronium. The annihilation usually results in two gamma rays. In most cases this direct annihilation occurs only after the positron has lost its excess kinetic energy and has thermalized with the material.
~10% of positrons form para-positronium, which then promptly (in ~0.12 ns) decays, usually into two gamma rays.
~30% of positrons form ortho-positronium but then annihilate within a few nanoseconds by 'picking off' another nearby electron with opposing spin. This usually produces two gamma rays. During this time, the very lightweight positronium atom exhibits a strong zero-point motion, that exerts a pressure and is able to push out a tiny nanometer-sized bubble in the medium.
Only ~0.5% of positrons form ortho-positronium that self-decays (usually into three gamma rays). This natural decay rate of ortho-positronium is relatively slow (~140 ns decay lifetime), compared to the aforementioned pick-off process, which is why the three-gamma decay rarely occurs.
History
The Croatian physicist Stjepan Mohorovičić predicted the existence of positronium in a 1934 article published in Astronomische Nachrichten, in which he called it the "electrum". Other sources incorrectly credit Carl Anderson as having predicted its existence in 1932 while at Caltech. It was experimentally discovered by Martin Deutsch at MIT in 1951 and became known as positronium. Many subsequent experiments have precisely measured its properties and verified predictions of quantum electrodynamics.
A discrepancy known as the ortho-positronium lifetime puzzle persisted for some time, but was resolved with further calculations and measurements. Measurements were in error because of the lifetime measurement of unthermalised positronium, which was produced at only a small rate. This had yielded lifetimes that were too long. Also calculations using relativistic quantum electrodynamics are difficult, so they had been done to only the first order. Corrections that involved higher orders were then calculated in a non-relativistic quantum electrodynamics.
In 2024, the AEgIS collaboration at CERN was the first to cool positronium by laser light, leaving it available for experimental use. The substance was brought to using laser cooling.
Exotic compounds
Molecular bonding was predicted for positronium. Molecules of positronium hydride (PsH) can be made. Positronium can also form a cyanide and can form bonds with halogens or lithium.
The first observation of di-positronium () molecules—molecules consisting of two positronium atoms—was reported on 12 September 2007 by David Cassidy and Allen Mills from University of California, Riverside.
Unlike muonium, positronium does not have a nucleus analogue, because the electron and the positron have equal masses. Consequently, while muonium tends to behave like a light isotope of hydrogen, positronium shows large differences in size, polarisability, and binding energy from hydrogen.
Natural occurrence
The events in the early universe leading to baryon asymmetry predate the formation of atoms (including exotic varieties such as positronium) by around a third of a million years, so no positronium atoms occurred then.
Likewise, the naturally occurring positrons in the present day result from high-energy interactions such as in cosmic ray–atmosphere interactions, and so are too hot (thermally energetic) to form electrical bonds before annihilation.
| Physical sciences | Atomic physics | Physics |
44599 | https://en.wikipedia.org/wiki/Chalcedony | Chalcedony | Chalcedony ( or ) is a cryptocrystalline form of silica, composed of very fine intergrowths of quartz and moganite. These are both silica minerals, but they differ in that quartz has a trigonal crystal structure, while moganite is monoclinic. Chalcedony's standard chemical structure (based on the chemical structure of quartz) is SiO2 (silicon dioxide).
Chalcedony has a waxy luster, and may be semitransparent or translucent. It can assume a wide range of colors, but those most commonly seen are white to gray, grayish-blue or a shade of brown ranging from pale to nearly black. The color of chalcedony sold commercially is often enhanced by dyeing or heating.
The name chalcedony comes from the Latin (alternatively spelled ) and is probably derived from the town of Chalcedon in Turkey. The name appears in Pliny the Elder's as a term for a translucent kind of jaspis. Another reference to a gem by the name of () is found in the Book of Revelation (21:19); however, it is a hapax legomenon, found nowhere else in the Bible, so it is hard to tell whether the precious gem mentioned in Revelation is the same as the mineral known by this name today. The term plasma is sometimes used to refer to green translucent chalcedony.
Varieties
Chalcedony occurs in a wide range of varieties. Many semi-precious gemstones are in fact forms of chalcedony. The more notable varieties of chalcedony are as follows:
Agate
Agate is a variety of chalcedony characterized by either transparency or color patterns, such as multi-colored curved or angular banding. Opaque varieties are sometimes referred to as jasper. Fire agate shows iridescent phenomena on a brown background; iris agate shows exceptional iridescence when light (especially pinpointed light) is shone through the stone. Landscape agate is chalcedony with a number of different mineral impurities making the stone resemble landscapes.
Carnelian
Carnelian (also spelled cornelian) is a clear-to-translucent reddish-brown variety of chalcedony. Its hue may vary from a pale orange to an intense almost-black coloration. Similar to carnelian is sard, which is brown rather than red.
Chrysoprase
Chrysoprase (also spelled chrysophrase) is a green variety of chalcedony, which has been colored by nickel oxide. (The darker varieties of chrysoprase are also referred to as prase. However, the term prase is also used to describe green quartz and to a certain extent is a color-descriptor, rather than a rigorously defined mineral variety.)
Blue-colored chalcedony is sometimes referred to as "blue chrysoprase" if the color is sufficiently rich, though it derives its color from the presence of copper and is largely unrelated to nickel-bearing chrysoprase.
Fire agate
Fire agate is a variety of chalcedony with inclusions of goethite or limonite causing an iridescent effect. It can display a wide range of iridescent colors including red, orange, yellow, green, blue, and purple.
Heliotrope
Heliotrope is a green variety of chalcedony, containing red inclusions of iron oxide that resemble drops of blood, giving heliotrope its alternative name of bloodstone. In a similar variety known as plasma, the spots are yellow instead.
Moss agate
Moss agate contains green filament-like inclusions, giving it the superficial appearance of moss or blue cheese. There is also tree agate which is similar to moss agate except it is solid white with green filaments whereas moss agate usually has a transparent background, so the "moss" appears in 3D. It is not a true form of agate, as it lacks agate's defining feature of concentric banding.
Chrome chalcedony
Chrome chalcedony is a green variety of chalcedony, which is colored by chromium compounds. It is also known as "mtorolite" when found in Zimbabwe and "chiquitanita" when found in Bolivia.
Onyx
Onyx is a variant of agate with black and white banding. Similarly, agate with brown, orange, red and white banding is known as sardonyx.
Chalcedony ice-blue
In Greenland, white to greyish chalcedony is known from volcanic strata of the Paleocene, in the Disko-Nuussuaq area (West Greenland) and from the Scoresby Sound area (East Greenland). A light blue variety of chalcedony is known from Illorsuit, formed in the volcanic rocks along the southern coast of the island. Because of its bluish, ice-like colour, it has the local name chalcedony "ice-blue".
History
Chalcedony was used in tool making as early as 32,000 BP in Central Australia where archaeological studies at sites in the Cleland Hills uncovered flakes from stone brought in from quarries many kilometres away. Pre-contact uses described in the twentieth century included ceremonial stone knives.
Chalcedony was used for green and yellow color in prehistoric cave paintings, for example at the Bhimbetka rock shelters. The chalcedony was ground to powder form then mixed with water and animal fat or tree resin or gum.
In the Bronze Age chalcedony was in use in the Mediterranean region; for example, on Minoan Crete at the Palace of Knossos, chalcedony seals have been recovered dating to circa 1800 BC. People living along the Central Asian trade routes used various forms of chalcedony, including carnelian, to carve intaglios, ring bezels (the upper faceted portion of a gem projecting from the ring setting), and beads that show strong Greco-Roman influence.
Fine examples of first century objects made from chalcedony, possibly Kushan, were found in recent years at Tillya-tepe in north-western Afghanistan. Hot wax would not stick to it so it was often used to make seal impressions.
The term chalcedony is derived from the name of the ancient Greek town Chalkedon in Asia Minor, in modern English usually spelled Chalcedon, today the Kadıköy district of Istanbul.
According to tradition, at least three varieties of chalcedony were used in the Jewish High Priest's Breastplate. (Jewish tradition states that Moses' brother Aaron wore the Breastplate, with inscribed gems representing the twelve tribes of Israel.) The Breastplate supposedly included jasper, chrysoprase and sardonyx, and there is some debate as to whether other agates were also used.
In the 19th century, Idar-Oberstein, Germany, became the world's largest chalcedony processing center, working mostly on agates. Most of these agates were from Latin America, in particular Brazil. Originally the agate carving industry around Idar and Oberstein was driven by local deposits that were mined in the 15th century. Several factors contributed to the re-emergence of Idar-Oberstein as agate center of the world: ships brought agate nodules back as ballast, thus providing extremely cheap transport. In addition, cheap labor and a superior knowledge of chemistry allowed them to dye the agates in any color with processes that were kept secret. Each mill in Idar-Oberstein had four or five grindstones. These were of red sandstone, obtained from Zweibrücken; and two men ordinarily worked together at the same stone.
Geochemistry
Structure
Chalcedony was once thought to be a fibrous variety of cryptocrystalline quartz. More recently however, it has been shown to also contain a monoclinic polymorph of quartz, known as moganite. The fraction, by mass, of moganite within a typical chalcedony sample may vary from less than 5% to over 20%. The existence of moganite was once regarded as dubious, but it is now officially recognised by the International Mineralogical Association.
Solubility
Chalcedony is more soluble than quartz under low-temperature conditions, despite the two minerals being chemically identical. Possible reasons include the existence of the moganite component, defects caused by Brazil twinning, and small crystal size.
Solubility of quartz and chalcedony in pure water
This table gives equilibrium concentrations of total dissolved silicon as calculated by PHREEQC (PH REdox EQuilibrium (in C language, USGS)) using the llnl.dat database.
| Physical sciences | Silicate minerals | Earth science |
44600 | https://en.wikipedia.org/wiki/Carnelian | Carnelian | Carnelian (also spelled cornelian) is a brownish-red mineral commonly used as a semiprecious stone. Similar to carnelian is sard, which is generally harder and darker; the difference is not rigidly defined, and the two names are often used interchangeably. Both carnelian and sard are varieties of the silica mineral chalcedony colored by impurities of iron oxide. The color can vary greatly, ranging from pale orange to an intense almost-black coloration. Significant localities include Yanacodo (Peru); Ratnapura (Sri Lanka); and Thailand. It has been found in Indonesia, Brazil, India, Russia (Siberia), and Germany. In the United States, the official State Gem of Maryland is also a variety of carnelian called Patuxent River stone.
History
The red variety of chalcedony has been known to be used as beads since the Early Neolithic in Bulgaria. The first faceted (with constant 16+16=32 facets on each side of the bead) carnelian beads are described from the Varna Chalcolithic necropolis (middle of the 5th millennium BC). The bow drill was used to drill holes into carnelian in Mehrgarh in the 4th–5th millennium BC.
Carnelian was recovered from Bronze Age Minoan layers at Knossos on Crete in a form that demonstrated its use in decorative arts; this use dates to approximately 1800 BC. Carnelian was used widely during Roman times to make engraved gems for signet or seal rings for imprinting a seal with wax on correspondence or other important documents, as hot wax does not stick to carnelian. Sard was used for Assyrian cylinder seals, Egyptian and Phoenician scarabs, and early Greek and Etruscan gems. The Hebrew odem (also translated as sardius), was the first stone in the High Priest's breastplate, a red stone, probably sard but perhaps red jasper. In Revelation 4:3, the One seated on the heavenly throne seen in the vision of John the apostle is said to "look like jasper and (sardius transliterated)." And likewise it is in Revelation 21:20 as one of the precious stones in the foundations of the wall of the heavenly city.
There is a Neo-Assyrian seal made of carnelian in the Western Asiatic Seals collection of the British Museum that shows Ishtar-Gula as a star goddess. She is holding a ring of royal authority and is seated on a throne. She is shown with the spade of Marduk (his symbol), Sibbiti (seven) gods, the stylus of Nabu and a worshiper. An 8th century BC carnelian seal from the collection of the Ashmolean Museum in Oxford shows Ishtar-Gula with her dog facing the spade of Marduk and his red dragon.
Etymology
Although now the more common term, "carnelian" is a 16th-century corruption of the 14th-century word "cornelian" (and its associated orthographies corneline and cornalyn). Cornelian, cognate with similar words in several Romance languages, comes from the Mediaeval Latin , itself derived from the Latin word , the cornel cherry, whose translucent red fruits resemble the stone. The Oxford English Dictionary calls "carnelian" a perversion of "cornelian," by subsequent analogy with the Latin word ("flesh"). According to Pliny the Elder, sard derived its name from the city of Sardis in Lydia from which it came, and according to others, may ultimately be related to the Persian word (sered, "yellowish-red"). Another possible derivation is from the Greek σάρξ (sarx, "flesh"); compare the surer etymology of onyx, which comes from Greek ὄνυξ (onyx, "claw, fingernail"), presumably because onyx with flesh-colored and white bands can resemble a fingernail.
Distinction between carnelian and sard
The names carnelian and sard are often used interchangeably, but they can also be used to describe distinct subvarieties. The general differences are as follows:
All of these properties vary across a continuum, so the boundary between carnelian and sard is inherently blurry.
| Physical sciences | Silicate minerals | Earth science |
44603 | https://en.wikipedia.org/wiki/Calcite | Calcite | Calcite is a carbonate mineral and the most stable polymorph of calcium carbonate (CaCO3). It is a very common mineral, particularly as a component of limestone. Calcite defines hardness 3 on the Mohs scale of mineral hardness, based on scratch hardness comparison. Large calcite crystals are used in optical equipment, and limestone composed mostly of calcite has numerous uses.
Other polymorphs of calcium carbonate are the minerals aragonite and vaterite. Aragonite will change to calcite over timescales of days or less at temperatures exceeding 300 °C, and vaterite is even less stable.
Etymology
Calcite is derived from the German , a term from the 19th century that came from the Latin word for lime, (genitive ) with the suffix -ite used to name minerals. It is thus a doublet of the word chalk.
When applied by archaeologists and stone trade professionals, the term alabaster is used not just as in geology and mineralogy, where it is reserved for a variety of gypsum; but also for a similar-looking, translucent variety of fine-grained banded deposit of calcite.
Unit cell and Miller indices
In publications, two different sets of Miller indices are used to describe directions in hexagonal and rhombohedral crystals, including calcite crystals: three Miller indices in the directions, or four Bravais–Miller indices in the directions, where is redundant but useful in visualizing permutation symmetries.
To add to the complications, there are also two definitions of unit cell for calcite. One, an older "morphological" unit cell, was inferred by measuring angles between faces of crystals, typically with a goniometer, and looking for the smallest numbers that fit. Later, a "structural" unit cell was determined using X-ray crystallography. The morphological unit cell is rhombohedral, having approximate dimensions and , while the structural unit cell is hexagonal (i.e. a rhombic prism), having approximate dimensions and . For the same orientation, must be multiplied by 4 to convert from morphological to structural units. As an example, calcite cleavage is given as "perfect on {1 0 1}" in morphological coordinates and "perfect on {1 0 4}" in structural units. In indices, these are {1 0 1} and {1 0 4}, respectively. Twinning, cleavage and crystal forms are often given in morphological units.
Properties
The diagnostic properties of calcite include a defining Mohs hardness of 3, a specific gravity of 2.71 and, in crystalline varieties, a vitreous luster. Color is white or none, though shades of gray, red, orange, yellow, green, blue, violet, brown, or even black can occur when the mineral is charged with impurities.
Crystal habits
Calcite has numerous habits, representing combinations of over 1000 crystallographic forms. Most common are scalenohedra, with faces in the hexagonal directions (morphological unit cell) or {2 1 4} directions (structural unit cell); and rhombohedral, with faces in the or directions (the most common cleavage plane). Habits include acute to obtuse rhombohedra, tabular habits, prisms, or various scalenohedra. Calcite exhibits several twinning types that add to the observed habits. It may occur as fibrous, granular, lamellar, or compact. A fibrous, efflorescent habit is known as lublinite. Cleavage is usually in three directions parallel to the rhombohedron form. Its fracture is conchoidal, but difficult to obtain.
Scalenohedral faces are chiral and come in pairs with mirror-image symmetry; their growth can be influenced by interaction with chiral biomolecules such as L- and D-amino acids. Rhombohedral faces are not chiral.
Optical
Calcite is transparent to opaque and may occasionally show phosphorescence or fluorescence. A transparent variety called "Iceland spar" is used for optical purposes. Acute scalenohedral crystals are sometimes referred to as "dogtooth spar" while the rhombohedral form is sometimes referred to as "nailhead spar". The rhombohedral form may also have been the "sunstone" whose use by Viking navigators is mentioned in the Icelandic Sagas.
Single calcite crystals display an optical property called birefringence (double refraction). This strong birefringence causes objects viewed through a clear piece of calcite to appear doubled. The birefringent effect (using calcite) was first described by the Danish scientist Rasmus Bartholin in 1669. At a wavelength of about 590 nm, calcite has ordinary and extraordinary refractive indices of 1.658 and 1.486, respectively. Between 190 and 1700 nm, the ordinary refractive index varies roughly between 1.9 and 1.5, while the extraordinary refractive index varies between 1.6 and 1.4.
Thermoluminescence
Calcite has thermoluminescent properties mainly due to manganese divalent (). An experiment was conducted by adding activators such as ions of Mn, Fe, Co, Ni, Cu, Zn, Ag, Pb, and Bi to the calcite samples to observe whether they emitted heat or light. The results showed that adding ions (, , , , , , , , ) did not react. However, a reaction occurred when both manganese and lead ions were present in calcite. By changing the temperature and observing the glow curve peaks, it was found that and acted as activators in the calcite lattice, but was much less efficient than .
Measuring mineral thermoluminescence experiments usually use x-rays or gamma-rays to activate the sample and record the changes in glowing curves at a temperature of 700–7500 K. Mineral thermoluminescence can form various glow curves of crystals under different conditions, such as temperature changes, because impurity ions or other crystal defects present in minerals supply luminescence centers and trapping levels. Observing these curve changes also can help infer geological correlation and age determination.
Chemical
Calcite, like most carbonates, dissolves in acids by the following reaction
The carbon dioxide released by this reaction produces a characteristic effervescence when a calcite sample is treated with an acid.
Due to its acidity, carbon dioxide has a slight solubilizing effect on calcite. The overall reaction is
If the amount of dissolved carbon dioxide drops, the reaction reverses to precipitate calcite. As a result, calcite can be either dissolved by groundwater or precipitated by groundwater, depending on such factors as the water temperature, pH, and dissolved ion concentrations. When conditions are right for precipitation, calcite forms mineral coatings that cement rock grains together and can fill fractures. When conditions are right for dissolution, the removal of calcite can dramatically increase the porosity and permeability of the rock, and if it continues for a long period of time, may result in the formation of caves. Continued dissolution of calcium carbonate-rich formations can lead to the expansion and eventual collapse of cave systems, resulting in various forms of karst topography.
Calcite exhibits an unusual characteristic called retrograde solubility: it is less soluble in water as the temperature increases. Calcite is also more soluble at higher pressures.
Pure calcite has the composition . However, the calcite in limestone often contains a few percent of magnesium. Calcite in limestone is divided into low-magnesium and high-magnesium calcite, with the dividing line placed at a composition of 4% magnesium. High-magnesium calcite retains the calcite mineral structure, which is distinct from that of dolomite, . Calcite can also contain small quantities of iron and manganese. Manganese may be responsible for the fluorescence of impure calcite, as may traces of organic compounds.
Distribution
Calcite is found all over the world, and its leading global distribution is as follows:
United States
Calcite is found in many different areas in the United States. One of the best examples is the Calcite Quarry in Michigan. The Calcite Quarry is the largest carbonate mine in the world and has been in use for more than 85 years. Large quantities of calcite can be mined from these sizeable open pit mines.
Canada
Calcite can also be found throughout Canada, such as in Thorold Quarry and Madawaska Mine, Ontario, Canada.
Mexico
Abundant calcite is mined in the Santa Eulalia mining district, Chihuahua, Mexico.
Iceland
Large quantities of calcite in Iceland are concentrated in the Helgustadir mine. The mine was once the primary mining location of "Iceland spar." However, it currently serves as a nature reserve, and calcite mining will not be allowed.
England
Calcite is found in parts of England, such as Alston Moor, Egremont, and Frizington, Cumbria.
Germany
St. Andreasberg, Harz Mountains, and Freiberg, Saxony can find calcite.
Use and applications
Ancient Egyptians carved many items out of calcite, relating it to their goddess Bast, whose name contributed to the term alabaster because of the close association. Many other cultures have used the material for similar carved objects and applications.
A transparent variety of calcite known as Iceland spar may have been used by Vikings for navigating on cloudy days. A very pure crystal of calcite can split a beam of sunlight into dual images, as the polarized light deviates slightly from the main beam. By observing the sky through the crystal and then rotating it so that the two images are of equal brightness, the rings of polarized light that surround the sun can be seen even under overcast skies. Identifying the sun's location would give seafarers a reference point for navigating on their lengthy sea voyages.
In World War II, high-grade optical calcite was used for gun sights, specifically in bomb sights and anti-aircraft weaponry. It was used as a polarizer (in Nicol prisms) before the invention of Polaroid plates and still finds use in optical instruments. Also, experiments have been conducted to use calcite for a cloak of invisibility.
Microbiologically precipitated calcite has a wide range of applications, such as soil remediation, soil stabilization and concrete repair. It also can be used for tailings management and is designed to promote sustainable development in the mining industry.
Calcite can help synthesize precipitated calcium carbonate (PCC) (mainly used in the paper industry) and increase carbonation. Furthermore, due to its particular crystal habit, such as rhombohedron, hexagonal prism, etc., it promotes the production of PCC with specific shapes and particle sizes.
Calcite, obtained from an 80 kg sample of Carrara marble, is used as the IAEA-603 isotopic standard in mass spectrometry for the calibration of δ18O and δ13C.
Calcite can be formed naturally or synthesized. However, artificial calcite is the preferred material to be used as a scaffold in bone tissue engineering due to its controllable and repeatable properties.
Calcite can be used to alleviate water pollution caused by the excessive growth of cyanobacteria. Lakes and rivers can lead to cyanobacteria blooms due to eutrophication, which pollutes water resources. Phosphorus (P) is the leading cause of excessive growth of cyanobacteria. As an active capping material, calcite can help reduce P release from sediments into the water, thus inhibiting cyanobacteria overgrowth.
Natural occurrence
Calcite is a common constituent of sedimentary rocks, limestone in particular, much of which is formed from the shells of dead marine organisms. Approximately 10% of sedimentary rock is limestone. It is the primary mineral in metamorphic marble. It also occurs in deposits from hot springs as a vein mineral; in caverns as stalactites and stalagmites; and in volcanic or mantle-derived rocks such as carbonatites, kimberlites, or rarely in peridotites.
Cacti contain Ca-oxalate biominerals. Their death releases these biominerals into the environment, which subsequently transform to calcite via a monohydrocalcite intermediate, sequestering carbon.
Calcite is often the primary constituent of the shells of marine organisms, such as plankton (such as coccoliths and planktic foraminifera), the hard parts of red algae, some sponges, brachiopods, echinoderms, some serpulids, most bryozoa, and parts of the shells of some bivalves (such as oysters and rudists). Calcite is found in spectacular form in the Snowy River Cave of New Mexico as mentioned above, where microorganisms are credited with natural formations. Trilobites, which became extinct a quarter billion years ago, had unique compound eyes that used clear calcite crystals to form the lenses. It also forms a substantial part of birds' eggshells, and the δC of the diet is reflected in the δC of the calcite of the shell.
The largest documented single crystal of calcite originated from Iceland, measured and and weighed about 250 tons. Classic samples have been produced at Madawaska Mine, near Bancroft, Ontario.
Bedding parallel veins of fibrous calcite, often referred to in quarrying parlance as beef, occur in dark organic rich mudstones and shales, these veins are formed by increasing fluid pressure during diagenesis.
Formation processes
Calcite formation can proceed by several pathways, from the classical terrace ledge kink model to the crystallization of poorly ordered precursor phases like amorphous calcium carbonate (ACC) via an Ostwald ripening process, or via the agglomeration of nanocrystals.
The crystallization of ACC can occur in two stages. First, the ACC nanoparticles rapidly dehydrate and crystallize to form individual particles of vaterite. Second, the vaterite transforms to calcite via a dissolution and reprecipitation mechanism, with the reaction rate controlled by the surface area of a calcite crystal. The second stage of the reaction is approximately 10 times slower.
However, crystallization of calcite has been observed to be dependent on the starting pH and concentration of magnesium in solution. A neutral starting pH during mixing promotes the direct transformation of ACC into calcite without a vaterite intermediate. But when ACC forms in a solution with a basic initial pH, the transformation to calcite occurs via metastable vaterite, following the pathway outlined above. Magnesium has a noteworthy effect on both the stability of ACC and its transformation to crystalline CaCO3, resulting in the formation of calcite directly from ACC, as this ion destabilizes the structure of vaterite.
Epitaxial overgrowths of calcite precipitated on weathered cleavage surfaces have morphologies that vary with the type of weathering the substrate experienced: growth on physically weathered surfaces has a shingled morphology due to Volmer-Weber growth, growth on chemically weathered surfaces has characteristics of Stranski-Krastanov growth, and growth on pristine cleavage surfaces has characteristics of Frank - van der Merwe growth. These differences are apparently due to the influence of surface roughness on layer coalescence dynamics.
Calcite may form in the subsurface in response to microorganism activity, such as sulfate-dependent anaerobic oxidation of methane, where methane is oxidized and sulfate is reduced, leading to precipitation of calcite and pyrite from the produced bicarbonate and sulfide. These processes can be traced by the specific carbon isotope composition of the calcites, which are extremely depleted in the 13C isotope, by as much as −125 per mil PDB (δ13C).
In Earth history
Calcite seas existed in Earth's history when the primary inorganic precipitate of calcium carbonate in marine waters was low-magnesium calcite (lmc), as opposed to the aragonite and high-magnesium calcite (hmc) precipitated today. Calcite seas alternated with aragonite seas over the Phanerozoic, being most prominent in the Ordovician and Jurassic periods. Lineages evolved to use whichever morph of calcium carbonate was favourable in the ocean at the time they became mineralised, and retained this mineralogy for the remainder of their evolutionary history. Petrographic evidence for these calcite sea conditions consists of calcitic ooids, lmc cements, hardgrounds, and rapid early seafloor aragonite dissolution. The evolution of marine organisms with calcium carbonate shells may have been affected by the calcite and aragonite sea cycle.
Calcite is one of the minerals that has been shown to catalyze an important biological reaction, the formose reaction, and may have had a role in the origin of life. Interaction of its chiral surfaces (see Form) with aspartic acid molecules results in a slight bias in chirality; this is one possible mechanism for the origin of homochirality in living cells.
Climate change
Climate change is exacerbating ocean acidification, possibly leading to lower natural calcite production. The oceans absorb large amounts of from fossil fuel emissions into the air. The total amount of artificial absorbed by the oceans is calculated to be 118 ± 19 Gt C. If a large amount of dissolves in the sea, it will cause the acidity of the seawater to increase, thereby affecting the pH value of the ocean. Calcifying organisms in the sea, such as molluscs foraminifera, crustaceans, echinoderms and corals, are susceptible to pH changes. Meanwhile, these calcifying organisms are also an essential source of calcite. As ocean acidification causes pH to drop, carbonate ion concentrations will decline, potentially reducing natural calcite production.
Gallery
| Physical sciences | Minerals | Earth science |
44633 | https://en.wikipedia.org/wiki/Ultramarine | Ultramarine | Ultramarine is a deep blue color pigment which was originally made by grinding lapis lazuli into a powder. Its lengthy grinding and washing process makes the natural pigment quite valuable—roughly ten times more expensive than the stone it comes from and as expensive as gold.
The name ultramarine comes from the Latin . The word means 'beyond the sea', as the pigment was imported by Italian traders during the 14th and 15th centuries from mines in Afghanistan. Much of the expansion of ultramarine can be attributed to Venice which historically was the port of entry for lapis lazuli in Europe.
Ultramarine was the finest and most expensive blue used by Renaissance painters. It was often used for the robes of the Virgin Mary and symbolized holiness and humility. It remained an extremely expensive pigment until a synthetic ultramarine was invented in 1826.
Ultramarine is a permanent pigment when under ideal preservation conditions. Otherwise, it is susceptible to discoloration and fading.
Structure
The pigment consists primarily of a zeolite-based mineral containing small amounts of polysulfides. It occurs in nature as a proximate component of lapis lazuli containing a blue cubic mineral called lazurite. In the Colour Index International, the pigment of ultramarine is identified as P. Blue 29 77007.
The major component of lazurite is a complex sulfur-containing sodium-silicate (Na8–10Al6Si6O24S2–4), which makes ultramarine the most complex of all mineral pigments. Some chloride is often present in the crystal lattice as well. The blue color of the pigment is due to the radical anion, which contains an unpaired electron.
Visual properties
The best samples of ultramarine are a uniform deep blue while other specimens are of paler color.
Particle size distribution has been found to vary among samples of ultramarine from various workshops. Numerous grinding techniques used by painters have resulted in different pigment/medium ratios and particle size distributions. The grinding and purification process results in pigment with particles of various geometries. Different grades of pigment may have been used for different areas in a painting, a characteristic that is sometimes used in art authentication.
Shades and variations
International Klein Blue (IKB) a deep blue hue first mixed by the French artist Yves Klein.
Electric
Electric ultramarine is the tone of ultramarine that is halfway between blue and violet on the RGB (HSV) color wheel, the expression of the HSV color space of the RGB color model.
Production
Natural production
Historically, lapis lazuli stone was mined in Afghanistan and shipped overseas to Europe.
A method to produce ultramarine from lapis lazuli was introduced and later described by Cennino Cennini in the 15th century. This process consisted of grinding the lapis lazuli mineral, mixing the ground material with melted wax, resins, and oils, wrapping the resulting mass in a cloth, and then kneading it in a dilute lye solution, a potassium carbonate solution prepared by combining wood ash with water. The blue lazurite particles collect at the bottom of the pot, while the colorless crystalline material and other impurities remain at the top. This process was performed at least three times, with each successive extraction generating a lower quality material. The final extraction, consisting largely of colorless material as well as a few blue particles, brings forth ultramarine ash which is prized as a glaze for its pale blue transparency. This extensive process was specific to ultramarine because the mineral it comes from has a combination of both blue and colorless pigments. If an artist were to simply grind and wash lapis lazuli, the resulting powder would be a greyish-blue color that lacks purity and depth of color since lapis lazuli contains a high proportion of colorless material.
Although the lapis lazuli stone itself is relatively inexpensive, the lengthy process of pulverizing, sifting, and washing to produce ultramarine makes the natural pigment quite valuable and roughly ten times more expensive than the stone it comes from. The high cost of the imported raw material and the long laborious process of extraction combined has been said to make high-quality ultramarine as expensive as gold.
Synthetic production
In 1990, an estimated 20,000 tons of ultramarine were produced industrially. The raw materials used in the manufacture of synthetic ultramarine are the following:
white kaolin,
anhydrous sodium sulfate (Na2SO4),
anhydrous sodium carbonate (Na2CO3),
powdered sulfur,
powdered charcoal or relatively ash-free coal, or colophony in lumps.
The preparation is typically made in steps:
The first part of the process takes place at 700 to 750 °C in a closed furnace, so that sulfur, carbon and organic substances give reducing conditions. This yields a yellow-green product sometimes used as a pigment.
In the second step, air or sulfur dioxide at 350 to 450 °C is used to oxidize sulfide in the intermediate product to S2 and Sn chromophore molecules, resulting in the blue (or purple, pink or red) pigment.
The mixture is heated in a kiln, sometimes in brick-sized amounts.
The resultant solids are then ground and washed, as is the case in any other insoluble pigment's manufacturing process; the chemical reaction produces large amounts of sulfur dioxide. (Flue-gas desulfurization is thus essential to its manufacture where SO2 pollution is regulated.)
Ultramarine poor in silica is obtained by fusing a mixture of soft clay, sodium sulfate, charcoal, sodium carbonate, and sulfur. The product is at first white, but soon turns green "green ultramarine" when it is mixed with sulfur and heated. The sulfur burns, and a fine blue pigment is obtained. Ultramarine rich in silica is generally obtained by heating a mixture of pure clay, very fine white sand, sulfur, and charcoal in a muffle furnace. A blue product is obtained at once, but a red tinge often results. The different ultramarines—green, blue, red, and violet—are finely ground and washed with water.
Synthetic ultramarine is a more vivid blue than natural ultramarine, since the particles in synthetic ultramarine are smaller and more uniform than the particles in natural ultramarine and therefore diffuse light more evenly. Its color is unaffected by light nor by contact with oil or lime as used in painting. Hydrochloric acid immediately bleaches it with liberation of hydrogen sulfide. Even a small addition of zinc oxide to the reddish varieties especially causes a considerable diminution in the intensity of the color. Modern, synthetic ultramarine blue is a non-toxic, soft pigment that does not need much mulling to disperse into a paint formulation.
Structure and classification
Ultramarine is the aluminosilicate zeolite with a sodalite structure. Sodalite consists of interconnected aluminosilicate cages. Some of these cages contain polysulfide () groups that are the chromophore (color centre). The negative charge on these ions is balanced by ions that also occupy these cages.
The chromophore is proposed to be or S4.
History
Antiquity and Middle Ages
The name derives from Middle Latin , literally "beyond the sea" because it was imported from Asia by sea. In the past, it has also been known as azzurrum ultramarine, , , , . The current terminology for ultramarine includes natural ultramarine (English), (French), (German), (Italian), and (Spanish). The first recorded use of ultramarine as a color name in English was in 1598.
The first noted use of lapis lazuli as a pigment can be seen in 6th and 7th-century paintings in Zoroastrian and Buddhist cave temples in Afghanistan, near the most famous source of the mineral. Lapis lazuli has been identified in Chinese paintings from the 10th and 11th centuries, in Indian mural paintings from the 11th, 12th, and 17th centuries, and on Anglo-Saxon and Norman illuminated manuscripts from .
Ancient Egyptians used lapis lazuli in solid form for ornamental applications in jewelry, however, there is no record of them successfully formulating lapis lazuli into paint. Archaeological evidence and early literature reveal that lapis lazuli was used as a semi-precious stone and decorative building stone from early Egyptian times. The mineral is described by the classical authors Theophrastus and Pliny. There is no evidence that lapis lazuli was used ground as a painting pigment by ancient Greeks and Romans. Like ancient Egyptians, they had access to a satisfactory blue colorant in the synthetic copper silicate pigment, Egyptian blue.
Renaissance
Venice was central to both the manufacturing and distribution of ultramarine during the early modern period. The pigment was imported by Italian traders during the 14th and 15th centuries from mines in Afghanistan. Other European countries employed the pigment less extensively than in Italy; the pigment was not used even by wealthy painters in Spain at that time.
During the Renaissance, ultramarine was the finest and most expensive blue that could be used by painters. Color infrared photogenic studies of ultramarine in 13th and 14th-century Sienese panel paintings have revealed that historically, ultramarine has been diluted with white lead pigment in an effort to use the color more sparingly given its high price. The 15th century artist Cennino Cennini wrote in his painters' handbook: "Ultramarine blue is a glorious, lovely and absolutely perfect pigment beyond all the pigments. It would not be possible to say anything about or do anything to it which would not make it more so." Natural ultramarine is a difficult pigment to grind by hand, and for all except the highest quality of mineral, sheer grinding and washing produces only a pale grayish blue powder.
The pigment was most extensively used during the 14th through 15th centuries, as its brilliance complemented the vermilion and gold of illuminated manuscripts and Italian panel paintings. It was valued chiefly on account of its brilliancy of tone and its inertness in opposition to sunlight, oil, and slaked lime. It is, however, extremely susceptible to even minute and dilute mineral acids and acid vapors. Dilute HCl, HNO3, and H2SO4 rapidly destroy the blue color, producing hydrogen sulfide (H2S) in the process. Acetic acid attacks the pigment at a much slower rate than mineral acids.
Ultramarine was only used for frescoes when it was applied secco because frescoes' absorption rate made its use cost prohibitive. The pigment was mixed with a binding medium like egg to form a tempera and applied over dry plaster, such as in Giotto di Bondone's frescos in the Cappella degli Scrovegni or the Arena Chapel in Padua.
European artists used the pigment sparingly, reserving their highest quality blues for the robes of Mary and the Christ child, possibly in an effort to show piety, spending as a means of expressing devotion. As a result of the high price, artists sometimes economized by using a cheaper blue, azurite, for under painting. Most likely imported to Europe through Venice, the pigment was seldom seen in German art or art from countries north of Italy. Due to a shortage of azurite in the late 16th and 17th century, the price for the already-expensive ultramarine increased dramatically.
17th and 18th centuries
Johannes Vermeer made extensive use of ultramarine in his paintings. The turban of the Girl with a Pearl Earring is painted with a mixture of ultramarine and lead white, with a thin glaze of pure ultramarine over it. In Lady Standing at a Virginal, the young woman's dress is painted with a mixture of ultramarine and green earth, and ultramarine was used to add shadows in the flesh tones. Scientific analysis by the National Gallery in London of Lady Standing at a Virginal showed that the ultramarine in the blue seat cushion in the foreground had degraded and become paler with time; it would have been a deeper blue when originally painted.
19th century (invention of synthetic ultramarine)
The beginning of the development of artificial ultramarine blue is known from Goethe. In about 1787, he observed the blue deposits on the walls of lime kilns near Palermo in Sicily. He was aware of the use of these glassy deposits as a substitute for lapis lazuli in decorative applications. He did not mention if it was suitable to grind for a pigment.
In 1814, Tassaert observed the spontaneous formation of a blue compound, very similar to ultramarine, if not identical with it, in a lime kiln at St. Gobain. In 1824, this caused the to offer a prize for the artificial production of the precious color. Processes were devised by Jean Baptiste Guimet (1826) and by Christian Gmelin (1828), then professor of chemistry in Tübingen. While Guimet kept his process a secret, Gmelin published his, and became the originator of the "artificial ultramarine" industry.
Permanence
Easel paintings and illuminated manuscripts have revealed natural ultramarine in a perfect state of preservation even though the art may be several centuries old. In general, ultramarine is a permanent pigment. Although it is a sulfur-containing compound from which sulfur is readily emitted as H2S, historically, it has been mixed with lead white with no reported occurrences of the lead pigment blackening to become lead sulfide.
A plague known as "ultramarine sickness" has occasionally been observed among ultramarine oil paintings as a grayish or yellowish gray discoloration of the paint surface. This can occur with artificial ultramarine that is used industrially. The cause of this has been debated among experts, however, potential causes include atmospheric sulfur dioxide and moisture, acidity of an oil- or oleo-resinous paint medium, or slow drying of the oil during which time water may have been absorbed, creating swelling, opacity of the medium, and therefore whitening of the paint film.
Both natural and artificial ultramarine are stable to ammonia and caustic alkalis in ordinary conditions. Artificial ultramarine has been found to fade when in contact with lime when it is used to color concrete or plaster. These observations have led experts to speculate if the natural pigment's fading may be the result of contact with the lime plaster of fresco paintings.
Synthetic applications
Synthetic ultramarine, being very cheap, is used for wall painting, the printing of paper hangings, and calico. It also is used as a corrective for the yellowish tinge often present in things meant to be white, such as linen and paper. Bluing or "laundry blue" is a suspension of synthetic ultramarine, or the chemically different Prussian blue, that is used for this purpose when washing white clothes. It is often found in makeup such as mascaras or eye shadows.
Large quantities are used in the manufacture of paper, and especially for producing a kind of pale blue writing paper which was popular in Britain. During World War I, the RAF painted the outer roundels with a color made from ultramarine blue. This became BS 108(381C) aircraft blue. It was replaced in the 1960s by a new color made on phthalocyanine blue, called BS110(381C) roundel blue.
Terminology
Ultramarine is a blue made from natural lapis lazuli, or its synthetic equivalent which is sometimes called "French Ultramarine". More generally "ultramarine blue" can refer to a vivid blue.
The term ultramarine can also refer to other pigments. Variants of the pigment such as "ultramarine red," "ultramarine green," and "ultramarine violet" all resemble ultramarine with respect to their chemistry and crystal structure.
The term "ultramarine green" indicates a dark green while barium chromate is sometimes referred to as "ultramarine yellow". Ultramarine pigment has also been termed "Gmelin's Blue," "Guimet's Blue," "New blue," "Oriental Blue," and "Permanent Blue".
| Physical sciences | Colors | Physics |
44652 | https://en.wikipedia.org/wiki/Lazurite | Lazurite | Lazurite, old name Azure spar is a tectosilicate mineral with sulfate, sulfur and chloride with formula . It is a feldspathoid and a member of the sodalite group. Lazurite crystallizes in the isometric system although well‐formed crystals are rare. It is usually massive and forms the bulk of the gemstone lapis lazuli.
Mineral
Lazurite is a deep‐blue to greenish‐blue. The colour is due to the presence of anions. It has a Mohs hardness of 5.0 to 5.5 and a specific gravity of 2.4. It is translucent with a refractive index of 1.50. It is fusible at 3.5 on Wolfgang Franz von Kobell's fusibility scale, and soluble in HCl. It commonly contains or is associated with grains of pyrite.
Lazurite is a product of contact metamorphism of limestone and is typically associated with calcite, pyrite, diopside, humite, forsterite, hauyne and muscovite.
Other blue minerals, such as the carbonate mineral, azurite, and the phosphate mineral, lazulite, may be confused with lazurite, but are easily distinguished with careful examination. At one time, lazurite was a synonym for azurite.
Lazurite was first described in 1890 for an occurrence in the Sar-e-Sang District, Koksha Valley, Badakhshan Province, Afghanistan. It has been mined for more than 6,000 years in the lapis lazuli district of Badakhshan. It has been used as a pigment in painting and cloth dyeing since at least the 6th or 7th century. It is also mined at Lake Baikal in Siberia; Mount Vesuvius; Burma; Canada; and the United States. The name is from the Persian for blue.
The most important mineral component of lapis lazuli is lazurite (25% to 40%)
Redefinition
Most lapis lazuli gets its blue color from Hauyne and almost none contain "true lazurite". This was changed in 2021, as lazurite was redefined so that it is enough for a quarter (instead of half) of the cages to contain sulfide.
Structure
Lazurite and hauyne seem to have the same structure and both are sulfate-dominant minerals. Lazurite is a pigment (opalescent) and has a bright blue streak (especially as a component of the semiprecious stone lapis lazuli). Many hauynes have a white or pale blue streak and are translucent. The difference might be a consequence of the redox state (sulfate to sulfide ratio).
| Physical sciences | Silicate minerals | Earth science |
44668 | https://en.wikipedia.org/wiki/Porcupine | Porcupine | Porcupines are large rodents with coats of sharp spines, or quills, that protect them against predation. The term covers two families of animals: the Old World porcupines of the family Hystricidae, and the New World porcupines of the family Erethizontidae. Both families belong to the infraorder Hystricognathi within the profoundly diverse order Rodentia and display superficially similar coats of rigid or semi-rigid quills, which are modified hairs composed of keratin. Despite this, the two groups are distinct from one another and are not closely related to each other within the Hystricognathi. The largest species of porcupine is the third-largest living rodent in the world, after the capybara and beaver.
The Old World porcupines (Hystricidae) live in Italy, Asia (western and southern), and most of Africa. They are large, terrestrial, and strictly nocturnal.
The New World porcupines (Erethizontidae) are indigenous to North America and northern South America. They live in wooded areas and can climb trees, where some species spend their entire lives. They are less strictly nocturnal than their Old World counterparts and generally smaller.
Most porcupines are about long, with a long tail. Weighing , they are rounded, large, and slow, and use an aposematic strategy of defence. Porcupines' colouration consists of various shades of brown, grey and white. Porcupines' spiny protection resembles that of the only distantly related erinaceomorph hedgehogs and Australian monotreme echidnas as well as tenrecid tenrecs.
Etymology
The word porcupine comes from the Latin + , from Old Italian porcospino, . A regional American name for the animal is quill-pig.
A baby porcupine is a porcupette. When born, a porcupette's quills are soft hair; they harden within a few days, forming the sharp quills of adults.
Evolution
Fossils belonging to the genus Hystrix date back to the late Miocene of the continent of Africa.
Species
Taxonomy
A porcupine is any of 30 species of rodents belonging to the families Erethizontidae (genera: Coendou, Erethizon, and Chaetomys) or Hystricidae (genera: Atherurus, Hystrix, and Trichys). Porcupines vary in size considerably: Rothschild's porcupine of South America weighs less than a kilogram (2.2 lb); the crested porcupine found in Italy, North Africa, and sub-Saharan Africa can grow to well over . The two families of porcupines are quite different, and although both belong to the Hystricognathi branch of the vast order Rodentia, they are not closely related.
Old World compared with New World species
The 11 Old World porcupines tend to be fairly large and have spines grouped in clusters.
The two subfamilies of New World porcupines are mostly smaller (although the North American porcupine reaches about in length and ), have their quills attached singly rather than grouped in clusters, and are excellent climbers, spending much of their time in trees. The New World porcupines evolved their spines independently (through convergent evolution) and are more closely related to several other families of rodents than they are to the Old World porcupines.
Longevity
Porcupines have a relatively high longevity and hold the record for being the longest-living rodent, with one individual named Cooper living over 32 years.
Diet
The North American porcupine is a herbivore and often climbs trees for food; it eats leaves, herbs, twigs, and green plants such as clover. In the winter, it may eat bark. The African porcupine is not a climber; instead, it forages on the ground. It is mostly nocturnal but will sometimes forage for food during the day, eating bark, roots, fruits, berries, and farm crops. Porcupines have become a pest in Kenya and are eaten as a delicacy.
Defence
Defensive behaviour displays in a porcupine depend on sight, scent, and sound. Often, these displays are shown when a porcupine becomes agitated or annoyed. There are four main displays seen in a porcupine: (in order from least to most aggressive) quill erection, teeth clattering, odour emission, and attack. A porcupine's colouring aids in part of its defence as most of the predators are nocturnal and colour-blind. A porcupine's markings are black and white. The dark body and coarse hair of the porcupine are dark brown/black and when quills are raised, present a white strip down its back mimicking the look of a skunk. This, along with the raising of the sharp quills, deters predators. Along with the raising of the quills, porcupines clatter their teeth to warn predators not to approach. The incisors vibrate against each other, the strike zone shifts back, and the cheek teeth clatter. This behaviour is often paired with body shivering, which is used to further display the dangerous quills. The rattling of quills is aided by the hollow quills at the back end of the porcupine. The use of odour is when the sight and sound have failed. An unpleasant scent is produced from the skin above the tail in times of stress and is often seen with a quill erection. If these processes fail, the porcupine will attack by running sideways or backwards into predators. A porcupine's tail can also be swung in the direction of the predator; if contact is made, the quills could be impaled into the predator causing injury or death.
Quills
Porcupines' quills, or spines, take on various forms, depending on the species, but all are modified hairs coated with thick plates of keratin, and embedded in the skin musculature. Old World porcupines have quills embedded in clusters, whereas in New World porcupines, single quills are interspersed with bristles, underfur, and hair.
Quills are released by contact or may drop out when the porcupine shakes its body. New quills grow to replace lost ones. Despite what is commonly believed, porcupines cannot launch their quills at range.
There are some possible antibiotic properties within the quills, specifically associated with the free fatty acids coating the quills. The antibiotic properties are believed to aid a porcupine that has suffered from self-injury.
Uses by humans
Porcupines are seldom eaten in Western culture but are eaten often in Southeast Asia, particularly Vietnam, where the prominent use of them as a food source has contributed to declines in porcupine populations.
Naturalist William J. Long reported the taste of the North American porcupine as "vile" and "malodorous" and delightful only to a lover of strong cheese.
More commonly, their quills and guard hairs are used for traditional decorative clothing; for example, their guard hairs are used in the creation of the Native American "porky roach" headdress. The main quills may be dyed and then applied in combination with thread to embellish leather accessories, such as knife sheaths and leather bags. Lakota women would harvest the quills for quillwork by throwing a blanket over a porcupine and retrieving the quills left stuck in the blanket.
The presence of barbs, acting like anchors, causes increased pain when removing a quill that has pierced the skin. The shape of the barbs makes the quills effective for penetrating the skin and for remaining in place. The quills have inspired research for such applications as the design of hypodermic needles and surgical staples. In contrast to the current design for surgical staples, the porcupine quill and barb design would allow easy and painless insertion, as the staple would stay in the skin using the anchored barb design rather than being bent under the skin like traditional staples.
Porcupines are also sometimes kept as exotic pets.
Habitat
Porcupines occupy a small range of habitats in tropical and temperate parts of Asia, Southern Europe, Africa, and North and South America. They live in forests and deserts, rocky outcrops, and hillsides. Some New World porcupines live in trees, but Old World porcupines prefer a rocky environment. Porcupines can be found on rocky areas up to high. They are generally nocturnal but are occasionally active during daylight.
Classification
Porcupines are distributed into two evolutionarily independent groups within the suborder Hystricomorpha of the Rodentia.
Infraorder Hystricognathi
Family Hystricidae: Old World porcupines
African brush-tailed porcupine, Atherurus africanus
African crested porcupine, Hystrix cristata
Asiatic brush-tailed porcupine, Atherurus macrourus
Cape porcupine, Hystrix africaeaustralis
Indian porcupine, Hystrix indicus
Malayan porcupine, Hystrix brachyura
Himalayan porcupine, Hystrix (brachyura) hodgsoni
Sunda porcupine, Hystrix javanica
Sumatran porcupine, Hystrix (Thecurus) sumatrae
Thick-spined porcupine, Hystrix (Thecurus) crassispinis
Philippine porcupine, Hystrix (Thecurus) pumilis
Long-tailed porcupine, Trichys fasciculata
Parvorder Phiomorpha sensu stricto
Family Thryonomyidae: cane rats
Family Petromuridae: Dassie rats
Family Bathyergidae: African mole-rats
Parvorder Caviomorpha
Superfamily Erethizontoidea
Family Erethizontidae: New World porcupines
North American porcupine, Erethizon dorsatum
Brazilian porcupine, Coendou prehensilis
Bicolored-spined porcupine, Coendou bicolor
Andean porcupine, Coendou quichua
Black dwarf (Koopman's) porcupine, Coendou nycthemera (koopmani)
Rothschild's porcupine, Coendou rothschildi
Santa Marta porcupine, Coendou sanctemartae
Mexican hairy dwarf porcupine, Coendou mexicanus
Paraguaian hairy dwarf porcupine, Coendou spinosus
Bahia porcupine, Coendou insidiosus
Brown hairy dwarf porcupine, Coendou vestitus
Streaked dwarf porcupine, Coendou ichillus
Black-tailed hairy dwarf porcupine, Coendou melanurus
Roosmalen's dwarf porcupine, Coendou roosmalenorum
Frosted hairy dwarf porcupine, Coendou pruinosus
Stump-tailed porcupine, Coendou rufescens
Bristle-spined porcupine, Chaetomys subspinosus (sometimes considered an echimyid)
Superfamily Cavioidea
Family Hydrochaeridae: capybara
Family Caviidae: Guinea-pigs
Family Dasyproctidae: agoutis and acouchis
Superfamily Octodontoidea
Family Abrocomidae: chinchilla-rats
Family Octodontidae: degus
Family Ctenomyidae: tuco-tucos
Family Echimyidae: spiny rats
Family Myocastoridae: nutrias
Family Capromyidae: hutias
Superfamily Chinchilloidea
Family Chinchillidae: chinchillas and allies
Family Dinomyidae: pacaranas
| Biology and health sciences | Rodents | null |
44682 | https://en.wikipedia.org/wiki/CMYK%20color%20model | CMYK color model | The CMYK color model (also known as process color, or four color) is a subtractive color model, based on the CMY color model, used in color printing, and is also used to describe the printing process itself. The abbreviation CMYK refers to the four ink plates used: cyan, magenta, yellow, and key (most often black).
The CMYK model works by partially or entirely masking colors on a lighter, usually white, background. The ink reduces the light that would otherwise be reflected. Such a model is called subtractive because inks subtract some colors from white light; in the CMY model, white light minus red leaves cyan, white light minus green leaves magenta, and white light minus blue leaves yellow.
In additive color models, such as RGB, white is the additive combination of all primary colored lights, and black is the absence of light. In the CMYK model, it is the opposite: white is the natural color of the paper or other background, and black results from a full combination of colored inks. To save cost on ink, and to produce deeper black tones, unsaturated and dark colors are produced by using black ink instead of or in addition to combinations of cyan, magenta, and yellow.
The CMYK printing process was invented in the 1890s, when newspapers began to publish color comic strips.
Halftoning
With CMYK printing, halftoning (also called screening) allows for less than full saturation of the primary colors; tiny dots of each primary color are printed in a pattern small enough that humans perceive a solid color. Magenta printed with a 20% halftone, for example, produces a pink color, because the eye perceives the tiny magenta dots on the large white paper as lighter and less saturated than the color of pure magenta ink. Halftoning allows for a continuous variability of each color, which enables continuous color mixing of the primaries. Without halftoning, each primary would be binary, i.e. on/off, which only allows for the reproduction of eight colors: white, the three primaries (cyan, magenta, yellow), the three secondaries (red, green, blue), and black.
Comparison to CMY
The CMYK color model is based on the CMY color model, which omits the black ink. Four-color printing uses black ink in addition to subtractive primaries for several reasons:
In traditional preparation of color separations, a red keyline on the black line art marked the outline of solid or tint color areas. In some cases a black keyline was used when it served as both a color indicator and an outline to be printed in black because usually the black plate contained the keyline. The K in CMYK represents the keyline, or black, plate, also sometimes called the key plate.
Text is typically printed in black and includes fine detail (such as serifs). To avoid even slight blurring when reproducing text (or other finely detailed outlines) using three inks would require impractically accurate registration.
A combination of 100% cyan, magenta, and yellow inks soaks the paper with ink, making it slower to dry, causing bleeding, or (especially on low-quality paper such as newsprint) weakening the paper so much that it tears.
Although a combination of 100% cyan, magenta, and yellow inks would, in theory, completely absorb the entire visible spectrum of light and produce a perfect black, practical inks fall short of their ideal characteristics, and the result is a dark, muddy color that is not quite black. Black ink absorbs more light and yields much better blacks.
Black ink is less expensive than the combination of colored inks that makes black.
A black made with just CMY inks is sometimes called a composite black.
When a very dark area is wanted, a colored or gray CMY "bedding" is applied first, then a full black layer is applied on top, making a rich, deep black; this is called rich black.
The amount of black to use to replace amounts of the other inks is variable, and the choice depends on the technology, paper and ink in use. Processes called under color removal, under color addition, and gray component replacement are used to decide on the final mix; different CMYK recipes will be used depending on the printing task.
Other printer color models
CMYK, as well as all other process color printing, is contrasted with spot color printing, in which specific colored inks are used to generate the colors seen. Some printing presses are capable of printing with both four-color process inks and additional spot color inks at the same time. High-quality printed materials, such as marketing brochures and books, often include photographs requiring process-color printing, other graphic effects requiring spot colors (such as metallic inks), and finishes such as varnish, which enhances the glossy appearance of the printed piece.
CMYK are the process printers which often have a relatively small color gamut. Processes such as Pantone's proprietary six-color (CMYKOG) Hexachrome considerably expand the gamut. Light, saturated colors often cannot be created with CMYK, and light colors in general may make visible the halftone pattern. Using a CcMmYK process, with the addition of light cyan and magenta inks to CMYK, can solve these problems, and such a process is used by many inkjet printers, including desktop models.
Comparison with RGB displays
Comparisons between RGB displays and CMYK prints can be difficult, since the color reproduction technologies and properties are very different. A computer monitor mixes shades of red, green, and blue light to create color images. A CMYK printer instead uses light-absorbing cyan, magenta, and yellow inks, whose colors are mixed using dithering, halftoning, or some other optical technique.
Similar to electronic displays, the inks used in printing produce color gamuts that are only a subsets of the visible spectrum, and the two color modes have their own specific ranges, each being capable of producing colors the other is not. As a result of this, an image rendered on an electronic display and rendered in print can vary in appearance. When designing images to be printed, designers work in RGB color spaces (electronic displays) capable of rendering colors a CMYK process cannot, and it is often difficult to accurately visualize a printed result that must fit into a different color space that both lacks some colors an electronic display can produce and includes colors it cannot.
Spectrum of printed paper
To reproduce color, the CMYK color model codes for absorbing light rather than emitting it (as is assumed by RGB). The K component ideally absorbs all wavelengths and is therefore achromatic. The cyan, magenta, and yellow components are used for color reproduction and they may be viewed as the inverse of RGB: Cyan absorbs red, magenta absorbs green, and yellow absorbs blue (−R,−G,−B).
Conversion
Since RGB and CMYK spaces are both device-dependent spaces, there is no simple or general conversion formula that converts between them. Conversions are generally done through color management systems, using color profiles that describe the spaces being converted. An ICC profile defines the bidirectional conversion between a neutral "profile connection" color space (CIE XYZ or Lab) and a selected colorspace, in this case both RGB and CMYK. The precision of the conversion depends on the profile itself, the exact methodology, and because the gamuts do not generally match, the rendering intent and constraints such as ink limit.
ICC profiles, internally built out of lookup tables and other transformation functions, are capable of handling many effects of ink blending. One example is the dot gain, which show up as non-linear components in the color-to-density mapping. More complex interactions such as Neugebauer blending can be modelled in higher-dimension lookup tables.
The problem of computing a colorimetric estimate of the color that results from printing various combinations of ink has been addressed by many scientists. A general method that has emerged for the case of halftone printing is to treat each tiny overlap of color dots as one of 8 (combinations of CMY) or of 16 (combinations of CMYK) colors, which in this context are known as Neugebauer primaries. The resultant color would be an area-weighted colorimetric combination of these primary colors, except that the Yule–Nielsen effect of scattered light between and within the areas complicates the physics and the analysis; empirical formulas for such analysis have been developed, in terms of detailed dye combination absorption spectra and empirical parameters.
Standardization of printing practices allow for some profiles to be predefined. One of them is the US Specifications for Web Offset Publications, which has its ICC color profile built into some software including Microsoft Office (as Agfa RSWOP.icm).
| Physical sciences | Basics_7 | null |
44700 | https://en.wikipedia.org/wiki/Leprosy | Leprosy | Leprosy, also known as Hansen's disease (HD), is a long-term infection by the bacteria Mycobacterium leprae or Mycobacterium lepromatosis. Infection can lead to damage of the nerves, respiratory tract, skin, and eyes. This nerve damage may result in a lack of ability to feel pain, which can lead to the loss of parts of a person's extremities from repeated injuries or infection through unnoticed wounds. An infected person may also experience muscle weakness and poor eyesight. Leprosy symptoms may begin within one year, but for some people symptoms may take 20 years or more to occur.
Leprosy is spread between people, although extensive contact is necessary. Leprosy has a low pathogenicity, and 95% of people who contract or who are exposed to M. leprae do not develop the disease. Spread is likely through a cough or contact with fluid from the nose of a person infected by leprosy. Genetic factors and immune function play a role in how easily a person catches the disease. Leprosy does not spread during pregnancy to the unborn child or through sexual contact. Leprosy occurs more commonly among people living in poverty. There are two main types of the disease – paucibacillary and multibacillary, which differ in the number of bacteria present. A person with paucibacillary disease has five or fewer poorly pigmented, numb skin patches, while a person with multibacillary disease has more than five skin patches. The diagnosis is confirmed by finding acid-fast bacilli in a biopsy of the skin.
Leprosy is curable with multidrug therapy. Treatment of paucibacillary leprosy is with the medications dapsone, rifampicin, and clofazimine for six months. Treatment for multibacillary leprosy uses the same medications for 12 months. Several other antibiotics may also be used. These treatments are provided free of charge by the World Health Organization.
Leprosy is not highly contagious. People with leprosy can live with their families and go to school and work. In the 1980s, there were 5.2 million cases globally, but by 2020 this decreased to fewer than 200,000. Most new cases occur in one of 14 countries, with India accounting for more than half of all new cases. In the 20 years from 1994 to 2014, 16 million people worldwide were cured of leprosy. About 200 cases per year are reported in the United States. Central Florida accounted for 81% of cases in Florida and nearly 1 out of 5 leprosy cases nationwide. Separating people affected by leprosy by placing them in leper colonies is not supported by evidence but still occurs in some areas of India, China, Japan, Africa, and Thailand.
Leprosy has affected humanity for thousands of years. The disease takes its name from the Greek word (), from (; 'scale'), while the term "Hansen's disease" is named after the Norwegian physician Gerhard Armauer Hansen. Leprosy has historically been associated with social stigma, which continues to be a barrier to self-reporting and early treatment. Leprosy is classified as a neglected tropical disease. World Leprosy Day was started in 1954 to draw awareness to those affected by leprosy.
The study of leprosy and its treatment is known as leprology.
Signs and symptoms
Common symptoms present in the different types of leprosy include a runny nose; dry scalp; eye problems; skin lesions; muscle weakness; reddish skin; smooth, shiny, diffuse thickening of facial skin, ear, and hand; loss of sensation in fingers and toes; thickening of peripheral nerves; a flat nose from the destruction of nasal cartilages; and changes in phonation and other aspects of speech production. In addition, atrophy of the testes and impotence may occur.
Leprosy can affect people in different ways. The average incubation period is five years. People may begin to notice symptoms within the first year or up to 20 years after infection. The first noticeable sign of leprosy is often the development of pale or pink coloured patches of skin that may be insensitive to temperature or pain. Patches of discolored skin are sometimes accompanied or preceded by nerve problems including numbness or tenderness in the hands or feet. Secondary infections (additional bacterial or viral infections) can result in tissue loss, causing fingers and toes to become shortened and deformed, as cartilage is absorbed into the body. A person's immune response differs depending on the form of leprosy.
Approximately 30% of people affected with leprosy experience nerve damage. The nerve damage sustained is reversible when treated early but becomes permanent when appropriate treatment is delayed by several months. Damage to nerves may cause loss of muscle function, leading to paralysis. It may also lead to sensation abnormalities or numbness, which may lead to additional infections, ulcerations, and joint deformities.
Cause
M. leprae and M. lepromatosis
Mycobacterioum leprae and Mycobacterium lepromatosis are the mycobacteria that cause leprosy. M. lepromatosis is a relatively newly identified mycobacterium isolated from a fatal case of diffuse lepromatous leprosy in 2008. M. lepromatosis is indistinguishable clinically from M. leprae. M. leprae is an aerobic, rod-shaped, acid-fast bacterium with a waxy cell envelope characteristic of the genus Mycobacterium. M. leprae and M. lepromatosis are obligate intracellular pathogens and cannot grow or be cultured outside of host tissues. However, they can be grown using research animals such as mice and armadillos.
Naturally occurring infections have been reported in nonhuman primates (including the African chimpanzee, the sooty mangabey, and the cynomolgus macaque), armadillos, and red squirrels. Multilocus sequence typing of the armadillo M. leprae strains suggests that they were of human origin for at most a few hundred years. Thus, it is suspected that armadillos first acquired the organism incidentally from early European explorers of the Americas. This incidental transmission was sustained in the armadillo population, and it may be transmitted back to humans, making leprosy a zoonotic disease (spread between humans and animals).
Red squirrels (Sciurus vulgaris), a threatened species in Great Britain, were found to carry leprosy in November 2016. It has been suggested that the trade in red squirrel fur, highly prized in the medieval period and intensively traded, may have been responsible for the leprosy epidemic in medieval Europe. A pre-Norman era skull excavated in Hoxne, Suffolk, in 2017 was found to carry DNA from a strain of M. leprae which closely matched the strain carried by modern red squirrels on Brownsea Island.
Risk factors
The greatest risk factor for developing leprosy is contact with another person infected by leprosy. People who are exposed to a person who has leprosy are 5–8 times more likely to develop leprosy than members of the general population. Leprosy occurs more commonly among those living in poverty. Not all people who are infected with M. leprae develop symptoms.
Conditions that reduce immune function, such as malnutrition, other illnesses, or genetic mutations, may increase the risk of developing leprosy. Infection with HIV does not appear to increase the risk of developing leprosy. Certain genetic factors in the person exposed have been associated with developing lepromatous or tuberculoid leprosy.
Transmission
Transmission of leprosy occurs during close contact with those who are infected. Transmission of leprosy is through the upper respiratory tract. Older research suggested the skin as the main route of transmission, but research has increasingly favored the respiratory route. Transmission occurs through inhalation of bacilli present in upper airway secretion.
Leprosy is not sexually transmitted and is not spread through pregnancy to the unborn child. The majority (95%) of people who are exposed to M. leprae do not develop leprosy; casual contact such as shaking hands and sitting next to someone with leprosy does not lead to transmission. People are considered non-infectious 72 hours after starting appropriate multi-drug therapy. Two exit routes of M. leprae from the human body that are often described are the skin and the nasal mucosa, although their relative importance is not clear. Lepromatous cases show large numbers of organisms deep in the dermis, but whether they reach the skin surface in sufficient numbers is doubtful. Leprosy may also be transmitted to humans by armadillos, although the mechanism is not fully understood.
Genetics
{| class="wikitable" style="float: right; margin-left:15px; text-align:center"
|-
! Name
! Locus
! OMIM
! Gene
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| LPRS1
| 10p13
|
|
|-
| LPRS2
| 6q25
|
| PARK2, PACRG
|-'
| LPRS3
| 4q32
|
| TLR2|-
| LPRS4
| 6p21.3
|
| LTA|-
| LPRS5
| 4p14
|
| TLR1|-
| LPRS6
| 13q14.11
|
|
|}
Not all people infected or exposed to M. leprae develop leprosy, and genetic factors are suspected to play a role in susceptibility to an infection. Cases of leprosy often cluster in families, and several genetic variants have been identified. In many who are exposed, the immune system can eliminate the leprosy bacteria during the early infection stage before severe symptoms develop. A genetic defect in cell-mediated immunity may cause a person to be susceptible to develop leprosy symptoms after exposure to the bacteria. The region of DNA responsible for this variability is also involved in Parkinson's disease, giving rise to current speculation that the two disorders may be linked at the biochemical level.
Mechanism
Most leprosy complications are the result of nerve damage. The nerve damage occurs from direct invasion by the M. leprae bacteria and a person's immune response resulting in inflammation. The molecular mechanism underlying how M. leprae produces the symptoms of leprosy is not clear, but M. leprae has been shown to bind to Schwann cells, which may lead to nerve injury including demyelination and a loss of nerve function (specifically a loss of axonal conductance). Numerous molecular mechanisms have been associated with this nerve damage including the presence of a laminin-binding protein and the glycoconjugate (PGL-1) on the surface of M. leprae that can bind to laminin on peripheral nerves.
As part of the human immune response, white blood cell-derived macrophages may engulf M. leprae by phagocytosis. In the initial stages, small sensory and autonomic nerve fibers in the skin of a person with leprosy are damaged. This damage usually results in hair loss to the area, a loss of the ability to sweat, and numbness (decreased ability to detect sensations such as temperature and touch). Further peripheral nerve damage may result in skin dryness, more numbness, and muscle weaknesses or paralysis in the area affected. The skin can crack and if the skin injuries are not carefully cared for, there is a risk for a secondary infection that can lead to more severe damage.
Diagnosis
In countries where people are frequently infected, a person is considered to have leprosy if they have one of the following two signs:
Skin lesion consistent with leprosy and with definite sensory loss.
Positive skin smears.
Skin lesions can be single or many, and usually hypopigmented, although occasionally reddish or copper-colored. The lesions may be flat (macules), raised (papules), or solid elevated areas (nodular). Experiencing sensory loss at the skin lesion is a feature that can help determine if the lesion is caused by leprosy or by another disorder such as tinea versicolor. Thickened nerves are associated with leprosy and can be accompanied by loss of sensation or muscle weakness, but muscle weakness without the characteristic skin lesion and sensory loss is not considered a reliable sign of leprosy. In some cases, the presence of acid-fast leprosy bacilli in skin smears is considered diagnostic; however, the diagnosis is typically made without laboratory tests, based on symptoms. If a person has a new leprosy diagnosis and already has a visible disability caused by leprosy, the diagnosis is considered late.
In countries or areas where leprosy is uncommon, such as the United States, diagnosis of leprosy is often delayed because healthcare providers are unaware of leprosy and its symptoms. Early diagnosis and treatment prevent nerve involvement, the hallmark of leprosy, and the disability it causes.U.S. Department of Health and Human Services, Health Resources and Services Administration. (n.d.). National Hansen's disease (leprosy) program Retrieved from There is no recommended test to diagnose latent leprosy in people without symptoms. Few people with latent leprosy test positive for anti PGL-1. The presence of M. leprae bacterial DNA can be identified using a polymerase chain reaction (PCR)-based technique. This molecular test alone is not sufficient to diagnose a person, but this approach may be used to identify someone who is at high risk of developing or transmitting leprosy such as those with few lesions or an atypical clinical presentation. New approaches propose tools to diagnose leprosy through artificial intelligence.
Classification
Several different approaches for classifying leprosy exist. There are similarities between the classification approaches.
The World Health Organization (WHO) system distinguishes "paucibacillary" and "multibacillary" based on the proliferation of bacteria. ("pauci-" refers to a small quantity.)
The Ridley-Jopling scale provides five gradations.
The ICD-10, though developed by the WHO, uses Ridley-Jopling and not the WHO system. It also adds an indeterminate ("I") entry.
In MeSH, three groupings are used.
Leprosy may also occur with only neural involvement, without skin lesions.
Complications
Leprosy may cause the victim to lose limbs and digits but not directly. M. leprae attacks nerve endings and destroys the body's ability to feel pain and injury. Without feeling pain, people with leprosy have an increased risk of injuring themselves. Injuries become infected and result in tissue loss. Fingers, toes, and limbs become shortened and deformed as the tissue is absorbed into the body.
Prevention
Early disease detection is important, since physical and neurological damage may be irreversible even if cured. Medications can decrease the risk of those living with people who have leprosy from acquiring the disease and likely those with whom people with leprosy come into contact outside the home. The WHO recommends that preventive medicine be given to people who are in close contact with someone who has leprosy. The suggested preventive treatment is a single dose of rifampicin in adults and children over 2 years old who do not already have leprosy or tuberculosis. Preventive treatment is associated with a 57% reduction in infections within 2 years and a 30% reduction in infections within 6 years.
The Bacillus Calmette–Guérin (BCG) vaccine offers a variable amount of protection against leprosy in addition to its closely related target of tuberculosis. It appears to be 26% to 41% effective (based on controlled trials) and about 60% effective based on observational studies with two doses possibly working better than one. The WHO concluded in 2018 that the BCG vaccine at birth reduces leprosy risk and is recommended in countries with high incidence of TB and people who have leprosy. People living in the same home as a person with leprosy are suggested to take a BCG booster which may improve their immunity by 56%. Development of a more effective vaccine is ongoing.
A novel vaccine called LepVax entered clinical trials in 2017 with the first encouraging results reported on 24 participants published in 2020. If successful, this would be the first leprosy-specific vaccine available.
Treatment
Several leprostatic agents are available for treatment. A three-drug regimen of rifampicin, dapsone, and clofazimine is recommended for all people with leprosy, for six months for paucibacillary leprosy and 12 months for multibacillary leprosy.
Multidrug therapy (MDT) remains highly effective, and people are no longer infectious after the first monthly dose. MDT is safe and easy to use under field conditions because it is available in calendar-labelled blister packs. Post-treatment relapse rates remain low. Resistance has been reported in several countries, although the number of cases is small. People with rifampicin-resistant leprosy may be treated with second-line medications such as fluoroquinolones, minocycline, or clarithromycin, but the treatment duration is 24 months because of their lower bactericidal activity. Evidence on the potential benefits and harms of alternative regimens for drug-resistant leprosy is not available.
For people with nerve damage, protective footwear may help prevent ulcers and secondary infection. Canvas shoes may be better than PVC boots. There may be no difference between double rocker shoes and below-knee plaster. Topical ketanserin seems to have a better effect on ulcer healing than clioquinol cream or zinc paste, but the evidence for this is weak. Phenytoin applied to the skin improves skin changes to a greater degree when compared to saline dressings.
Outcomes
Although leprosy has been curable since the mid-20th century, left untreated it can cause permanent physical impairments and damage to a person's nerves, skin, eyes, and limbs. Despite leprosy not being very infectious and having a low pathogenicity, there is still significant stigma and prejudice associated with the disease. Because of this stigma, leprosy can affect a person's participation in social activities and may also affect the lives of their family and friends. People with leprosy are also at a higher risk for problems with their mental well-being. The social stigma may contribute to problems obtaining employment, financial difficulties, and social isolation. Efforts to reduce discrimination and reduce the stigma surrounding leprosy may help improve outcomes for people with leprosy.
Epidemiology
In 2018, there were 208,619 new cases of leprosy recorded, a slight decrease from 2017. In 2015, 94% of the new leprosy cases were confined to 14 countries. India reported the greatest number of new cases (60% of reported cases), followed by Brazil (13%) and Indonesia (8%). Although the number of cases worldwide continues to fall, there are parts of the world where leprosy is more common, including Brazil, South Asia (India, Nepal, Bhutan), some parts of Africa (Tanzania, Madagascar, Mozambique), and the western Pacific. About 150 to 250 cases are diagnosed in the United States each year.
In the 1960s, there were tens of millions of leprosy cases recorded when the bacteria started to develop resistance to dapsone, the most common treatment option at the time. International (e.g., the WHO's "Global Strategy for Reducing Disease Burden Due to Leprosy") and national (e.g., the International Federation of Anti-Leprosy Associations) initiatives have reduced the total number and the number of new cases of the disease.
The number of new leprosy cases is difficult to measure and monitor because of leprosy's long incubation period, delays in diagnosis after the onset of the disease, and lack of medical care in affected areas. The registered prevalence of the disease is used to determine disease burden. Registered prevalence is a useful proxy indicator of the disease burden, as it reflects the number of active leprosy cases diagnosed with the disease and receiving treatment with MDT at a given point in time. The prevalence rate is defined as the number of cases registered for MDT treatment among the population in which the cases have occurred, again at a given point in time.
History
Historical distribution
Using comparative genomics, in 2005, geneticists traced the origins and worldwide distribution of leprosy from East Africa or the Near East along human migration routes. They found four strains of M. leprae with specific regional locations: Monot et al. (2005) determined that leprosy originated in East Africa or the Near East and traveled with humans along their migration routes, including those of trade in goods and slaves. The four strains of M. leprae are based in specific geographic regions where each predominantly occurs:
strain 1 in Asia, the Pacific region, and East Africa;
strain 2 in Ethiopia, Malawi, Nepal, north India, and New Caledonia;
strain 3 in Europe, North Africa, and the Americas;
strain 4 in West Africa and the Caribbean.
This confirms the spread of the disease along the migration, colonisation, and slave trade routes taken from East Africa to India, West Africa to the New World, and from Africa to Europe and vice versa.
Skeletal remains discovered in 2009 represent the oldest documented evidence for leprosy, dating to the 2nd millennium BC. Located at Balathal, Rajasthan, in northwest India, the discoverers suggest that if the disease did migrate from Africa to India during the 3rd millennium BC "at a time when there was substantial interaction among the Indus Civilization, Mesopotamia, and Egypt, there needs to be additional skeletal and molecular evidence of leprosy in India and Africa to confirm the African origin of the disease". A proven human case was verified by DNA taken from the shrouded remains of a man discovered by researchers from the Hebrew University of Jerusalem in a tomb in Akeldama, next to the Old City of Jerusalem dated by radiocarbon methods to the first half of the 1st century.
The oldest strains of leprosy known from Europe are from Great Chesterford in southeast England and date back to AD 415–545. These findings suggest a different path for the spread of leprosy, meaning it may have originated in western Eurasia. This study also indicates that there were more strains in Europe at the time than previously determined.
Discovery and scientific progress
Literary attestation of leprosy is unclear because of the ambiguity of many early sources, including the Indian Atharvaveda and Kausika Sutra, the Egyptian Ebers Papyrus, and the Hebrew Bible's various sections regarding signs of impurity (tzaraath). Leprotic symptoms are attested in the Indian doctor Sushruta's Compendium, originally dating to c. 600 BC but only surviving in emended texts no earlier than the 5th century BC. Symptoms consistent with leprosy were possibly described by Hippocrates in 460 BC. However, Hansen's disease probably did not exist in Greece or the Middle East before the Common Era. In 1846, Francis Adams produced The Seven Books of Paulus Aegineta which included a commentary on all medical and surgical knowledge and descriptions and remedies to do with leprosy from the Romans, Greeks, and Arabs.Roman: Celsus, Pliny, Serenus Samonicus, Scribonius Largus, Caelius Aurelianus, Themison, Octavius Horatianus, Marcellus the Emperic; Greek: Aretaeus, Plutarch, Galen, Oribasius, Aetius (Aëtius of Amida or Sicamus Aëtius), Actuarius, Nonnus, Psellus, Leo, Myrepsus; Arabic: Scrapion, Avenzoar, Albucasis, Haly Abbas translated by Stephanus Antiochensis, Alsharavius, Rhases (Abū Bakr al-Rāzī), and Guido de Cauliaco.
Leprosy did not exist in the Americas before colonization by modern Europeans nor did it exist in Polynesia until the middle of the 19th century. The causative agent of leprosy, M. leprae, was discovered by Gerhard Armauer Hansen in Norway in 1873, making it one of the first species of pathogenic bacteria identified.
Treatment
Chaulmoogra tree oil was used topically to manage Hansen's disease for centuries. Chaulmoogra oil could not be taken orally without causing nausea or injected without forming an abscess. Leprosy was once believed to be highly contagious and was treated with mercury, as was syphilis, which was first described in 1530. Many early cases thought to be leprosy could actually have been syphilis. In 1915, Alice Ball, the first black woman to graduate from the University of Hawai'i with a masters in chemistry, discovered how to make the oil water soluble. This technique led to marked improvements in patients with Hansen's disease who were treated in Hawai'i.
The first effective drug (promin) became available in the 1940s. In the 1950s, dapsone was introduced. The search for further effective antileprosy drugs led to the use of clofazimine and rifampicin in the 1960s and 1970s. Later, Indian scientist Shantaram Yawalkar and his colleagues formulated a combined therapy using rifampicin and dapsone, intended to mitigate bacterial resistance. Combining all three drugs was first recommended by the WHO in 1981. These three drugs are still used in the standard MDT regimens. Resistance has developed to initial treatment. Until the introduction of MDT in the early 1980s, leprosy could not be diagnosed and treated successfully within the community.
The importance of the nasal mucosa in the transmission of M. leprae was recognized as early as 1898 by Schäffer, in particular, that of the ulcerated mucosa. The mechanism of plantar ulceration in leprosy and its treatment was first described by Ernest W. Price.
Etymology
The word "leprosy" comes from the Greek word "λέπος (lépos) – skin" and "λεπερός (leperós) – scaly man".
Society and culture
Treatment cost
Between 1995 and 1999 the WHO, with the aid of the Nippon Foundation, supplied all endemic countries with free MDT in blister packs, channeled through ministries of health. This free provision was extended in 2000 and again in 2005, 2010, and 2015 with donations by the MDT manufacturer Novartis through the WHO. At the national level, non-governmental organizations (NGOs) affiliated with the national program will continue to be provided with an appropriate free supply.
Historical texts
Written accounts of leprosy date back thousands of years. Various skin diseases translated as leprosy appear in the ancient Indian text, the Atharava Veda, by 600 BC. Another Indian text, the Manusmriti (200 BC), prohibit contact with those infected with the disease and make marriage to a person infected with leprosy punishable.
The Hebraic root tsara or tsaraath (צָרַע, – tsaw-rah' – to be struck with leprosy, to be leprous) and the Greek (λεπρός – lepros), are of broader classification than the more narrow use of the term related to Hansen's Disease. Any progressive skin disease (a whitening or splotchy bleaching of the skin, raised manifestations of scales, scabs, infections, rashes, etc.) — as well as generalized molds and surface discoloration of any clothing, leather, or discoloration on walls or surfaces throughout homes — all came under the "law of leprosy" (Leviticus 14:54–57). Ancient sources such as the Talmud (Sifra 63) make clear that tzaraath refers to various types of lesions or stains associated with ritual impurity and occurring on cloth, leather, or houses, as well as skin. Traditional Judaism and Jewish rabbinical authorities, both historical and modern, emphasize that the tsaraath of Leviticus is a spiritual ailment with no direct relationship to Hansen's disease or physical contagions. The relation of tsaraath to "leprosy" comes from translations of Hebrew Biblical texts into Greek and ensuing misconceptions.
All three Synoptic Gospels of the New Testament describe instances of Jesus healing people with leprosy (Matthew 8:1–4, Mark 1:40–45, and Luke 5:12–16). The Bible's description of leprosy is congruous (if lacking detail) with the symptoms of modern leprosy, but the relationship between this disease, tzaraath, and Hansen's disease has been disputed. The biblical perception that people with leprosy were unclean can be found in a passage from Leviticus 13:44–46. While this text defines the leper as impure, it does not explicitly make a moral judgement on those with leprosy. Some early Christians believed that those affected by leprosy were being punished by God for sinful behavior. Moral associations have persisted throughout history. In the 6th century Pope Gregory the Great and Isidore of Seville considered people with the disease to be heretics.
Middle Ages
The social perception of leprosy in the general population was mixed. On one hand, people feared getting infected with the disease and thought of people suspected of leprosy to be unclean, untrustworthy, and occasionally morally corrupt. On the other hand, Jesus' interaction with lepers, the writing of church leaders, and the Christian focus on charitable works led to viewing the lepers as "chosen by God" or seeing the disease as a means of obtaining access to heaven.
Early medieval understanding of leprosy was influenced by early Christian writers such as Gregory of Nazianzus and John Chrysostom, whose writings were later embraced by Byzantine and Latin writers. Gregory, for example, composed sermons urging Christians to assist victims of the disease, and he condemned pagans or Christians who justified rejecting lepers on the allegation that God had sent them the disease to punish them. As cases of leprosy increased in the Eastern Roman Empire, becoming a major health issue, the ecclesiastic leaders discussed how to assist those affected as well as how to change the attitude of society towards them. They also tried this by using the name "holy disease" instead of the commonly used "elephant's disease" (elephantiasis), implying that God did not create this disease to punish people but to purify them for heaven. Although not always successful in persuading the public and a cure was never found by Greek medicians, they created an environment where victims could get palliative care and were never expressly banned from society, as sometimes happened in western Europe. Theodore Balsamon, a 12th-century jurist in Constantinople, noted that lepers were allowed to enter the same churches, cities, and assemblies that healthy people attended.
As the disease became more prevalent in western Europe in the 5th century, efforts began to set up permanent institutions to house and feed lepers. These efforts were, inclusively, the work of bishops in France at the end of the sixth century, such as in Chalon-sur-Saône. The increase in hospitals or leprosaria (sing. leprosarium) that treated people with leprosy in the 12th and 13th century seems to indicate a rise in cases, possibly in connection with the increase in urbanisation as well as returning crusaders from the Middle East. France alone had nearly 2,000 leprosaria during this period. Additionally to the new leprosia, further steps were taken by secular and religious leaders to prevent further spread of the disease. The third Lateran Council of 1179 required lepers to have their own priests and churches and a 1346 edict by King Edward expelled lepers from city limits. Segregation from mainstream society became common, and people with leprosy were often required to wear clothing that identified them as such or carry a bell announcing their presence. As in the East, it was the Church who took care of the lepers due to the persisting moral stigma and who ran the leprosaria. Although the leprosaria in Western Europe removed the sick from society, they were never a place to quarantine them or from which they could not leave: lepers would go beg for alms for the upkeep of the leprosaria or meet with their families.
Multiple groups in Western Europe from the Middle Ages faced social ostracization and discrimination that was justified, in part, due to claims that they were the descendants of lepers. These groups included the Cagots and the Caquins.
19th century
Norway was the location of a progressive stance on leprosy tracking and treatment and played an influential role in European understanding of the disease. In 1832, Dr. JJ Hjort conducted the first leprosy survey, thus establishing a basis for epidemiological surveys. Subsequent surveys resulted in the establishment of a national leprosy registry to study the causes of leprosy and to track the rate of infection. Early leprosy research throughout Europe was conducted by Norwegian scientists Daniel Cornelius Danielssen and Carl Wilhelm Boeck. Their work resulted in the establishment of the National Leprosy Research and Treatment Center. Danielssen and Boeck believed the cause of leprosy transmission was hereditary. This stance was influential in advocating for the isolation of those infected by sex to prevent reproduction.
Though leprosy rates were on the decline in the Western world by the 1860s, authorities frequently embraced isolation treatment due to a combination of reasons, including fears of the disease spreading from the Global South, efforts by Christian missionaries and a lack of understanding concerning bacteriology, medical diagnosis and how contagious the disease was. The rapid expansion of Western imperialism during the Victorian era resulted in westerners coming into increasing contact with regions where the disease was endemic, including British India. English surgeon Henry Vandyke Carter observed isolation treatment for leprosy patients first-hand while visiting Norway, applying these methods in British India with the financial and logistical assistance of Protestant missionaries. Colonialist and religious viewpoints of the disease continued to be a major factor in the treatment and public perception of the disease in the Global South until decolonization in the mid-20th century.
20th and 21st century
In 1898, the colonial government in British India enacted the Leprosy Act of 1898, which mandated the compulsory segregation of people with leprosy by authorities in newly established leper asylums, where they were segregated by sex to prevent sexual activity. The act, which proved difficult to enforce, was repealed in 1983 by the Indian government after MDT had become widely available in India. In 1983, the National Leprosy Elimination Programme, previously the National Leprosy Control Programme, changed its methods from surveillance to the treatment of people with leprosy. India still accounts for over half of the global disease burden. According to WHO, new cases in India during 2019 diminished to 114,451 patients (57% of the world's total new cases). Until 2019, Indians could justify a petition for divorce with their spouse's diagnosis of leprosy.
The National Leprosarium at Carville, Louisiana, known in 1955 as the Louisiana Leper Home, was the only leprosy hospital in the mainland United States. Leprosy patients from all over the United States were sent to Carville to be kept in isolation away from the public, as not much about leprosy transmission was known at the time and stigma against those with leprosy was high. The Carville leprosarium was known for its innovations in reconstructive surgery for those with leprosy. In 1941, 22 patients at Carville underwent trials for a new drug called promin. The results were described as miraculous, and soon after the success of promin came dapsone, a medicine even more effective in the fight against leprosy. Leprosy incidence peaked in the United States in 1983, followed by a steep decline. However, case numbers have been slowly rising again since 2000. In 2020, 159 cases were reported in the country.
Stigma
Despite now effective treatment and education efforts, leprosy stigma continues to be problematic in developing countries where the disease is common. Leprosy is most common amongst impoverished populations where social stigma is likely to be compounded by poverty. Fears of ostracism, loss of employment, or expulsion from family and society may contribute to a delayed diagnosis and treatment.
Folk beliefs, lack of education, and religious connotations of the disease continue to influence social perceptions of those affected in many parts of the world. In Brazil, for example, folklore holds that leprosy is a disease transmitted by dogs, or that it is associated with sexual promiscuity, or that it is a punishment for sins or moral transgressions (distinct from other diseases and misfortunes, which are in general thought of as being according to the will of God). Socioeconomic factors also have a direct impact. Lower-class domestic workers who are often employed by those in a higher socioeconomic class may find their employment in jeopardy as physical manifestations of the disease become apparent. Skin discoloration and darker pigmentation resulting from the disease also have social repercussions.
In extreme cases in northern India, leprosy is equated with an "untouchable" status that "often persists long after individuals with leprosy have been cured of the disease, creating lifelong prospects of divorce, eviction, loss of employment, and ostracism from family and social networks."
Public policy
A goal of the WHO is to "eliminate leprosy," and in 2016 the organization launched "Global Leprosy Strategy 2016–2020: Accelerating towards a leprosy-free world". Elimination of leprosy is defined as "reducing the proportion of (people with) leprosy in the community to very low levels, specifically to below one case per 10,000 population". Diagnosis and treatment with multidrug therapy are effective, and a 45% decline in disease burden has occurred since MDT has become more widely available. The organization emphasizes the importance of fully integrating leprosy treatment into public health services, effective diagnosis and treatment, and access to information. The approach includes supporting an increase in health care professionals who understand the disease, and a coordinated and renewed political commitment that includes coordination between countries and improvements in the methodology for collecting and analysing data.
Interventions include:
Early detection of cases focusing on children to reduce transmission and disabilities.
Enhanced healthcare services and improved access for people who may be marginalized.
For countries where leprosy is endemic, further interventions include an improved screening of close contacts, improved treatment regimens, and interventions to reduce stigma and discrimination against people who have leprosy.
In some instances in India, community-based rehabilitation is embraced by local governments and NGOs alike. Often, the identity cultivated by a community environment is preferable to reintegration, and models of self-management and collective agency independent of NGOs and government support have been desirable and successful.
Notable cases
Josephine Cafrine of Seychelles (1877–1907) had leprosy from age 12 and kept a journal that documented her struggles and suffering. It was published as an autobiography in 1923.
Saint Damien De Veuster (1840–1889), a Catholic priest from Belgium, contracted leprosy and ministered to lepers who had been placed under a government-sanctioned medical quarantine on Molokai in the Kingdom of Hawaii.
Baldwin IV of Jerusalem (1161–1185), Catholic king of Latin Jerusalem, had leprosy.
Ingeborg Grytten ( – after 1705), Norwegian writer whose leprosy is thought to have influenced her poetry, which is characterized by a strong religious faith in God's salvation.
Josefina Guerrero (1917–1996), Filipino World War II spy who used the Japanese fear of her leprosy to listen to their battle plans and deliver the information to the American forces under Douglas MacArthur.
King Henry IV of England () possibly had leprosy.
Vietnamese poet Hàn Mặc Tử (1912–1940)
Ōtani Yoshitsugu (1558 or 1565–1600), a Japanese daimyō (feudal lord).
British writer Peter Greave (1910–1977).
In media
— the short story collection Life's Handicap by Rudyard Kipling has a story "The Mark of the Beast" in which a traveller on horseback literally stumbles into a leper colony in India.
— Jack London published "Koolau the Leper" in his Tales of Hawai'i about Molokai and people consigned to it circa 1893.
— James Michener's novel Hawaii dramatizes Molokai's leper settlement, including Father Damien.
— Ben-Hur depicts the title character's mother, Miriam, and younger sister, Tirzah, are imprisoned by the Roman Empire. When they are freed years later, they have leprosy and leave town for the Valley of the Lepers, rather than stay and reunite with Ben-Hur. They leave the colony, and when Jesus dies on the cross, they are miraculously cured.
— English author Graham Greene's novel A Burnt-Out Case is set in a leper colony in Belgian Congo. The story is also predominantly about a disillusioned architect working with a doctor on devising new cures and amenities for mutilated victims of lepers; the title, too, refers to the condition of mutilation and disfigurement in the disease.
— Forugh Farrokhzad made a 22-minute documentary about a leprosy colony in Iran titled The House Is Black. The film humanizes the people affected and opens by saying that "there is no shortage of ugliness in the world, but by closing our eyes on ugliness, we will intensify it."
— The lead character in The Chronicles of Thomas Covenant by Stephen R. Donaldson suffers from leprosy. His condition seems to be cured by the magic of the fantasy land he finds himself in, but he resists believing in its reality, for example, by continuing to perform a regular visual surveillance of extremities as a safety check. Donaldson gained experience with the disease as a young man in India, where his father worked in a missionary for people with leprosy.
— Moloka'i is a novel by Alan Brennert about a leper colony in Hawaii. This novel follows the story of a seven-year-old girl taken from her family and put on Molokai's leper settlement.
— Squint: My Journey with Leprosy is a memoir by Jose P. Ramirez.
Infection of animals
Between 15 and 20% of nine-banded armadillos (Dasypus novemcinctus) in the south-central United States carry M. leprae.Liang, Jiayu (10/16/2023). "Leprosy in Florida: medical experts monitoring unusual, new cases of Hansen's disease". University of Florida Emerging Pathogens Institute. https://epi.ufl.edu/2023/10/16/leprosy-in-florida-medical-experts-monitoring-unusual-new-cases-of-hansens-disease/#:~:text=A%20major%20hint%20that%20someone,leprae%20bacteria. As a result of their low body temperature their tissues commonly contain massive numbers of organisms which helps in the dissemination of the infection. Armadillos were first demonstrated in 1971 to develop leprosy after inoculation with M. leprae. Because of armadillos' armor, skin lesions are difficult to ascertain. Abrasions around the eyes, nose, and feet are the most common signs. Infected armadillos make up a large reservoir of M. leprae and may be a source of infection for some humans in the United States or other locations in the armadillos' home range. In armadillo leprosy, lesions do not persist at the site of entry in animals; M. leprae multiply in macrophages at the site of inoculation and lymph nodes.
Armadillos have been used in immunological research to fight leprosy. Some notable reagents include recombinant interleukin-2 and recombinant interferon-gamma reagents. Additionally, they have been key and have been useful models of leprosy for studies regarding neuropathy. In clinical procedures such as electrophysiological nerve conduction tests Armadillo's nerve function has been properly assessed. Despite the studies mentioned regarding Armadillo's relationship to neuropathy and other effects of leprosy, there is still a lack of proper study on armadillos, and in conducting more armadillo-specific regents our understanding of leprosy’s effects on armadillos and possible humans can be found. Armadillos are a key component of modern-day research on leprosy.
There is a stigma surrounding armadillos and the carrying of leprosy. Because many people do not understand armadillos very well it is common for people to think of them as being dangerous to society and as a result valuing their lives less than other animals. It has become more common in parts of America for people to eat raw or undercooked armadillo, making the chances high that if not properly handled with care, one may become infected.
An outbreak in chimpanzees in West Africa is showing that the bacteria can infect another species and also possibly have additional rodent hosts. Studies have demonstrated that the disease is endemic in the UK red Eurasian squirrel population, with M. leprae and M. lepromatosis appearing in different populations. The M. leprae strain discovered on Brownsea Island is equated to one thought to have died out in the human population in mediaeval times. Despite this, and speculation regarding past transmission through trade in squirrel furs, there does not seem to be a high risk of squirrel to human transmission from the wild population. Although leprosy continues to be diagnosed in immigrants to the UK, the last known human case of leprosy arising in the UK was recorded over 200 years ago.
It has been shown that leprosy can reprogram cells in mouse and armadillo models, similar to how induced pluripotent stem cells are generated by the transcription factors Myc, Oct3/4, Sox2, and Klf4. A notable study conducted by Charles Shepard used mice to find how leprosy, an infection that has a preference for cooler areas of the body, would work in a warm-blooded animal. This mice study helped further the understanding of how leprosy works in humans. This was called "The Mouse Model." The main findings were that even in mice whose immune systems were severely impaired and at a perceived high risk of developing leprosy, the body was still in most cases able to fight off leprosy. The findings suggest that in mice, the body will use their body's energy to fight leprosy. Using The Mouse Model Shepard was able to conduct new research regarding leprosy. This model can now be used as a tool to further study M. leprae. The Mouse Model takes a more easily accessible animal model to better understand this complex disease. There are a few other up-and-coming models for M. leprae'' including the use of other animals including but not limited to mammals, birds, and cold-blooded animals. These animals do not tend to give as great results as armadillos and mice as different animals have different levels of disease resistance.
| Biology and health sciences | Infectious disease | null |
44723 | https://en.wikipedia.org/wiki/Printing | Printing | Printing is a process for mass reproducing text and images using a master form or template. The earliest non-paper products involving printing include cylinder seals and objects such as the Cyrus Cylinder and the Cylinders of Nabonidus. The earliest known form of printing evolved from ink rubbings made on paper or cloth from texts on stone tablets, used during the sixth century. Printing by pressing an inked image onto paper (using woodblock printing) appeared later that century. Later developments in printing technology include the movable type invented by Bi Sheng around 1040 and the printing press invented by Johannes Gutenberg in the 15th century. The technology of printing played a key role in the development of the Renaissance and the Scientific Revolution and laid the material basis for the modern knowledge-based economy and the spread of learning to the masses.
History
Woodblock printing
Woodblock printing is a technique for printing text, images or patterns that was used widely throughout East Asia. It originated in China in antiquity as a method of printing on textiles and later on paper.
In East Asia
The earliest examples of ink-squeeze rubbings and potential stone printing blocks appear in the mid-sixth century in China. A type of printing called mechanical woodblock printing on paper started during the 7th century in the Tang dynasty, and subsequently spread throughout East Asia. Nara Japan printed the Hyakumantō Darani en masse around 770, and distributed them to temples throughout Japan. In Korea, an example of woodblock printing from the eighth century was discovered in 1966. A copy of the Buddhist Dharani Sutra called the Pure Light Dharani Sutra (), discovered in Gyeongju, in a Silla dynasty pagoda that was repaired in 751, was undated but must have been created sometime before the reconstruction of the Shakyamuni Pagoda of Bulguk Temple, Kyongju Province in 751. The document is estimated to have been created no later than 704.
By the ninth century, printing on paper had taken off, and the first completely surviving printed book is the Diamond Sutra (British Library) of 868, uncovered from Dunhuang. By the tenth century, 400,000 copies of some sutras and pictures were printed, and the Confucian classics were in print. A skilled printer could print up to 2,000 double-page sheets per day.
Printing spread early to Korea and Japan, which also used Chinese logograms, but the technique was also used in Turpan and Vietnam using a number of other scripts. This technique then spread to Persia and Russia. This technique was transmitted to Europe by around 1400 and was used on paper for old master prints and playing cards.
In the Middle East
Block printing, called tarsh in Arabic, developed in Arabic Egypt during the ninth and tenth centuries, mostly for prayers and amulets. There is some evidence to suggest that these print blocks were made from non-wood materials, possibly tin, lead, or clay. The techniques employed are uncertain. Block printing later went out of use during the Timurid Renaissance. The printing technique in Egypt was embraced by reproducing texts on paper strips and supplying them in different copies to meet the demand.
In Europe
Block printing first came to Europe as a method for printing on cloth, where it was common by 1300. Images printed on cloth for religious purposes could be quite large and elaborate. When paper became relatively easily available, around 1400, the technique transferred very quickly to small woodcut religious images and playing cards printed on paper. These prints were produced in very large numbers from about 1425 onward.
Around the mid-fifteenth-century, block-books, woodcut books with both text and images, usually carved in the same block, emerged as a cheaper alternative to manuscripts and books printed with movable type. These were all short, heavily illustrated works, the bestsellers of the day, repeated in many different block-book versions: the and the Biblia pauperum were the most common. There is still some controversy among scholars as to whether their introduction preceded or, in the majority view, followed the introduction of movable type, with the estimated range of dates being between about 1440 and 1460.
Movable-type printing
Movable type is the system of printing and typography using movable pieces of metal type, made by casting from matrices struck by letterpunches. Movable type allowed for much more flexible processes than hand copying or block printing.
Around 1040, the first known movable type system was created in China by Bi Sheng out of porcelain. Bi Sheng used clay type, which broke easily, but Wang Zhen by 1298 had carved a more durable type from wood. He also developed a complex system of revolving tables and number-association with written Chinese characters that made typesetting and printing more efficient. Still, the main method in use there remained woodblock printing (xylography), which "proved to be cheaper and more efficient for printing Chinese, with its thousands of characters".
Copper movable type printing originated in China at the beginning of the 12th century. It was used in large-scale printing of paper money
issued by the Northern Song dynasty. Movable type spread to Korea during the Goryeo dynasty.
Around 1230, Koreans invented a metal type movable printing using bronze. The Jikji, published in 1377, is the earliest known metal printed book. Type-casting was used, adapted from the method of casting coins. The character was cut in beech wood, which was then pressed into a soft clay to form a mould, and bronze poured into the mould, and finally the type was polished. Eastern metal movable type was spread to Europe between the late 14th and early 15th centuries. The Korean form of metal movable type was described by the French scholar Henri-Jean Martin as "extremely similar to Gutenberg's". Authoritative historians Frances Gies and Joseph Gies claimed that "The Asian priority of invention movable type is now firmly established, and that Chinese-Korean technique, or a report of it traveled westward is almost certain."
The printing press
Around 1450, Johannes Gutenberg introduced the first movable type printing system in Europe. He advanced innovations in casting type based on a matrix and hand mould, adaptations to the screw-press, the use of an oil-based ink, and the creation of a softer and more absorbent paper. Gutenberg was the first to create his type pieces from an alloy of lead, tin, antimony, copper and bismuth – the same components still used today. Johannes Gutenberg started work on his printing press around 1436, in partnership with Andreas Dritzehen – whom he had previously instructed in gem-cutting – and Andreas Heilmann, the owner of a paper mill.
Compared to woodblock printing, movable type page setting and printing using a press was faster and more durable. Also, the metal type pieces were sturdier and the lettering more uniform, leading to typography and fonts. The high quality and relatively low price of the Gutenberg Bible (1455) established the superiority of movable type for Western languages. The printing press rapidly spread across Europe, leading up to the Renaissance, and later all around the world.
Time Life magazine called Gutenberg's innovations in movable type printing the most important invention of the second millennium.
Rotary printing press
The steam-powered rotary printing press, invented in 1843 in the United States by Richard M. Hoe, ultimately allowed millions of copies of a page in a single day. Mass production of printed works flourished after the transition to rolled paper, as continuous feed allowed the presses to run at a much faster pace. Hoe's original design operated at up to 2,000 revolutions per hour where each revolution deposited 4 page images, giving the press a throughput of 8,000 pages per hour. By 1891, The New York World and Philadelphia Item were operating presses producing either 90,000 4-page sheets per hour or 48,000 8-page sheets.
The rotary printing press uses impressions curved around a cylinder to print on long continuous rolls of paper or other substrates. Rotary drum printing was later significantly improved by William Bullock. There are multiple types of rotary printing press technologies that are still used today: sheetfed offset, rotogravure, and flexographic printing.
Printing capacity
The table lists the maximum number of pages which various press designs could print per hour.
Conventional printing technology
All printing process are concerned with two kinds of areas on the final output:
Image area (printing areas)
Non-image area (non-printing areas)
After the information has been prepared for production (the prepress step), each printing process has definitive means of separating the image from the non-image areas.
Conventional printing has four types of process:
Planographics, in which the printing and non-printing areas are on the same plane surface and the difference between them is maintained chemically or by physical properties, the examples are: offset lithography, collotype, and screenless printing.
Relief, in which the printing areas are on a plane surface and the non printing areas are below the surface, examples: flexography and letterpress.
Intaglio, in which the non-printing areas are on a plane surface and the printing area are etched or engraved below the surface, examples: steel die engraving, gravure, etching, collagraph.
Porous or Stencil, in which the printing areas are on fine mesh screens through which ink can penetrate, and the non-printing areas are a stencil over the screen to block the flow of ink in those areas, examples: screen printing, stencil duplicator, risograph.
Crop marks
To print an image without a blank area around the image, the non-printing areas must be trimmed after printing. Crop marks can be used to show the printer where the printing area ends, and the non-printing area begins. The part of the image which is trimmed off is called bleed.
Letterpress
Letterpress printing is a technique of relief printing. A worker composes and locks movable type into the bed of a press, inks it, and presses paper against it to transfer the ink from the type which creates an impression on the paper.
There is different paper for different works the quality of paper shows different ink to use.
Letterpress printing was the normal form of printing text from its invention by Johannes Gutenberg in the mid-15th century and remained in wide use for books and other uses until the second half of the 20th century, when offset printing was developed. More recently, letterpress printing has seen a revival in an artisanal form.
Offset
Offset printing is a widely used modern printing process. This technology is best described as when a positive (right-reading) image on a printing plate is inked and transferred (or "offset") from the plate to a rubber blanket. The blanket image becomes a mirror image of the plate image. An offset transfer moves the image to a printing substrate (typically paper), making the image right-reading again. Offset printing uses a lithographic process which is based on the repulsion of oil and water. The offset process employs a flat (planographic) image carrier (plate) which is mounted on a press cylinder. The image to be printed obtains ink from ink rollers, while the non-printing area attracts an (acidic) film of water, keeping the non-image areas ink-free. Most offset presses use three cylinders: Plate, blanket, impression. Currently, most books and newspapers are printed using offset lithography.
Gravure
Gravure printing is an intaglio printing technique, where the image being printed is made up of small depressions in the surface of the printing plate. The cells are filled with ink, and the excess is scraped off the surface with a doctor blade. Then a rubber-covered roller presses paper onto the surface of the plate and into contact with the ink in the cells. The printing cylinders are usually made from copper plated steel, which is subsequently chromed, and may be produced by diamond engraving; etching, or laser ablation.
Gravure printing is known for its ability to produce high-quality, high-resolution images with accurate color reproduction and using viscosity control equipment during production.
Ink evaporation control affects the change in the color of the printed image.
Gravure printing is used for long, high-quality print runs such as magazines, mail-order catalogues, packaging and printing onto fabric and wallpaper. It is also used for printing postage stamps and decorative plastic laminates, such as kitchen worktops.
Flexography
Flexography is a type of relief printing. The relief plates are typically made from photopolymers.
The process is used for flexible packaging, corrugated board, labels, newspapers and more. In this market it competes with gravure printing by holding 80% of the market in US, 50% in Europe but only 20% in Asia.
Other printing techniques
The other significant printing techniques include:
Dye-sublimation printer
Inkjet, used typically to print a small number of books or packaging, and also to print a variety of materials: from high quality papers simulating offset printing, to floor tiles. Inkjet is also used to apply mailing addresses to direct mail pieces
Laser printing (toner printing) mainly used in offices and for transactional printing (bills, bank documents). Laser printing is commonly used by direct mail companies to create variable data letters or coupons.
Pad printing, popular for its ability to print on complex three-dimensional surfaces
Relief print, mainly used for catalogues
Screen printing for a variety of applications ranging from T-shirts to floor tiles, and on uneven surfaces
Intaglio, used mainly for high value documents such as currencies.
Thermal printing, popular in the 1990s for fax printing. Used today for printing labels such as airline baggage tags and individual price labels in supermarket deli counters.
Impact of German movable type printing press
Quantitative aspects
It is estimated that following the innovation of Gutenberg's printing press, the European book output rose from a few million to around one billion copies within a span of less than four centuries.
Religious impact
Samuel Hartlib, who was exiled in Britain and enthusiastic about social and cultural reforms, wrote in 1641 that "the art of printing will so spread knowledge that the common people, knowing their own rights and liberties, will not be governed by way of oppression".
In the Muslim world, printing, especially in Arabic scripts, was strongly opposed throughout the early modern period, partially due to the high artistic renown of the art of traditional calligraphy. However, printing in Hebrew or Armenian script was often permitted. Thus, the first movable type printing in the Ottoman Empire was in Hebrew in 1493, after which both religious and non-religious texts were able to be printed in Hebrew. According to an imperial ambassador to Istanbul in the middle of the sixteenth century, it was a sin for the Turks, particularly Turkish Muslims, to print religious books. In 1515, Sultan Selim I issued a decree under which the practice of printing would be punishable by death. At the end of the sixteenth century, Sultan Murad III permitted the sale of non-religious printed books in Arabic characters, yet the majority were imported from Italy. Ibrahim Muteferrika established the first press for printing in Arabic in the Ottoman Empire, against opposition from the calligraphers and parts of the Ulama. It operated until 1742, producing altogether seventeen works, all of which were concerned with non-religious, utilitarian matters. Printing did not become common in the Islamic world until the 19th century.
Hebrew language printers were banned from printing guilds in some Germanic states; as a result, Hebrew printing flourished in Italy, beginning in 1470 in Rome, then spreading to other cities including Bari, Pisa, Livorno, and Mantua. Local rulers had the authority to grant or revoke licenses to publish Hebrew books, and many of those printed during this period carry the words 'con licenza de superiori' (indicating their printing having been officially licensed) on their title pages.
It was thought that the introduction of printing 'would strengthen religion and enhance the power of monarchs.' The majority of books were of a religious nature, with the church and crown regulating the content. The consequences of printing 'wrong' material were extreme. Meyrowitz used the example of William Carter who in 1584 printed a pro-Catholic pamphlet in Protestant-dominated England. The consequence of his action was hanging.
Social impact
Print gave a broader range of readers access to knowledge and enabled later generations to build directly on the intellectual achievements of earlier ones without the changes arising within verbal traditions. Print, according to Acton in his 1895 lecture On the Study of History, gave "assurance that the work of the Renaissance would last, that what was written would be accessible to all, that such an occultation of knowledge and ideas as had depressed the Middle Ages would never recur, that not an idea would be lost".
Print was instrumental in changing the social nature of reading.
Elizabeth Eisenstein identifies two long-term effects of the invention of printing. She claims that print created a sustained and uniform reference for knowledge and allowed comparisons of incompatible views.
Asa Briggs and Peter Burke identify five kinds of reading that developed in relation to the introduction of print:
Critical reading: Because texts finally became accessible to the general population, critical reading emerged as people were able to form their own opinions on texts.
Dangerous reading: Reading was seen as a dangerous pursuit because it was considered rebellious and unsociable, especially in the case of women, because reading could stir up dangerous emotions such as love, and if women could read, they could read love notes.
Creative reading: Printing allowed people to read texts and interpret them creatively, often in very different ways than the author intended.
Extensive reading: Once print made a wide range of texts available, earlier habits of intensive reading of texts from start to finish began to change, and people began reading selected excerpts, allowing much more extensive reading on a wider range of topics.
Private reading: Reading was linked to the rise of individualism because, before print, reading was often a group event in which one person would read to a group. With print, both literacy and the availability of texts increased, and solitary reading became the norm.
The invention of printing also changed the occupational structure of European cities. Printers emerged as a new group of artisans for whom literacy was essential, while the much more labour-intensive occupation of the scribe naturally declined. Proof-correcting arose as a new occupation, while a rise in the numbers of booksellers and librarians naturally followed the explosion in the numbers of books.
Educational impact
Gutenberg's printing press had profound impacts on universities as well. Universities were influenced in their "language of scholarship, libraries, curriculum, [and] pedagogy"
The language of scholarship
Before the invention of the printing press, most written material was in Latin. However, after the invention of printing the number of books printed expanded as well as the vernacular. Latin was not replaced completely, but remained an international language until the eighteenth century.
University libraries
At this time, universities began establishing accompanying libraries. "Cambridge made the chaplain responsible for the library in the fifteenth century but this position was abolished in 1570 and in 1577 Cambridge established the new office of university librarian. Although, the University of Leuven did not see a need for a university library based on the idea that professor were the library. Libraries also began receiving so many books from gifts and purchases that they began to run out of room. However, the issue was solved in 1589 by a man named Merton who decided books should be stored on horizontal shelves rather than lecterns.
Curriculum
The printed press changed university libraries in many ways. Professors were finally able to compare the opinions of different authors rather than being forced to look at only one or two specific authors. Textbooks themselves were also being printed in different levels of difficulty, rather than just one introductory text being made available.
Comparison of printing methods
Digital printing
By 2005, digital printing accounted for approximately 9% of the 45 trillion pages printed annually around the world.
Printing at home, an office, or an engineering environment is subdivided into:
small format (up to ledger size paper sheets), as used in business offices and libraries
wide format (up to 3' or 914mm wide rolls of paper), as used in drafting and design establishments.
Some of the more common printing technologies are:
blueprint – and related chemical technologies
daisy wheel – where pre-formed characters are applied individually
dot-matrix – which produces arbitrary patterns of dots with an array of printing studs
line printing – where formed characters are applied to the paper by lines
heat transfer – such as early fax machines or modern receipt printers that apply heat to special paper, which turns black to form the printed image
inkjet – including bubble-jet, where ink is sprayed onto the paper to create the desired image
electrophotography – where toner is attracted to a charged image and then developed
laser – a type of xerography where the charged image is written pixel by pixel using a laser
solid ink printer – where solid sticks of ink are melted to make liquid ink or toner
Vendors typically stress the total cost to operate the equipment, involving complex calculations that include all cost factors involved in the operation as well as the capital equipment costs, amortization, etc. For the most part, toner systems are more economical than inkjet in the long run, even though inkjets are less expensive in the initial purchase price.
Professional digital printing (using toner) primarily uses an electrical charge to transfer toner or liquid ink to the substrate onto which it is printed. Digital print quality has steadily improved from early color and black and white copiers to sophisticated colour digital presses such as the Xerox iGen3, the Kodak Nexpress, the HP Indigo Digital Press series, and the InfoPrint 5000. The iGen3 and Nexpress use toner particles and the Indigo uses liquid ink. The InfoPrint 5000 is a full-color, continuous forms inkjet drop-on-demand printing system. All handle variable data, and rival offset in quality. Digital offset presses are also called direct imaging presses, although these presses can receive computer files and automatically turn them into print-ready plates, they cannot insert variable data.
Small press and fanzines generally use digital printing. Prior to the introduction of cheap photocopying, the use of machines such as the spirit duplicator, hectograph, and mimeograph was common.
3D printing
3D printing is a form of manufacturing technology where physical objects are created from three-dimensional digital models using 3D printers. The objects are created by laying down or building up many thin layers of material in succession. The technique is also known as additive manufacturing, rapid prototyping, or fabricating.
In the 1980s, 3D printing techniques were considered suitable only for the production of functional or aesthetic prototypes, and a more appropriate term for it at the time was rapid prototyping. , the precision, repeatability, and material range of 3D printing have increased to the point that some 3D printing processes are considered viable as an industrial-production technology, whereby the term additive manufacturing can be used synonymously with 3D printing. One of the key advantages of 3D printing is the ability to produce very complex shapes or geometries that would be otherwise infeasible to construct by hand, including hollow parts or parts with internal truss structures to reduce weight. Fused deposition modeling (FDM), which uses a continuous filament of a thermoplastic material, is the most common 3D printing process in use .
| Technology | Media and communication | null |
44731 | https://en.wikipedia.org/wiki/Calcium%20carbonate | Calcium carbonate | Calcium carbonate is a chemical compound with the chemical formula . It is a common substance found in rocks as the minerals calcite and aragonite, most notably in chalk and limestone, eggshells, gastropod shells, shellfish skeletons and pearls. Materials containing much calcium carbonate or resembling it are described as calcareous. Calcium carbonate is the active ingredient in agricultural lime and is produced when calcium ions in hard water react with carbonate ions to form limescale. It has medical use as a calcium supplement or as an antacid, but excessive consumption can be hazardous and cause hypercalcemia and digestive issues.
Chemistry
Calcium carbonate shares the typical properties of other carbonates. Notably it
reacts with acids, releasing carbonic acid which quickly disintegrates into carbon dioxide and water:
releases carbon dioxide upon heating, called a thermal decomposition reaction, or calcination (to above 840 °C in the case of ), to form calcium oxide, CaO, commonly called quicklime, with reaction enthalpy 178 kJ/mol:
reacts with gaseous hydrogen to form methane and water vapor plus solid calcium oxide or calcium hydroxide depending on temperature and product gas composition. Various metals including palladium and nickel are catalysts for the reaction.
Calcium carbonate reacts with water that is saturated with carbon dioxide to form the soluble calcium bicarbonate.
This reaction is important in the erosion of carbonate rock, forming caverns, and leads to hard water in many regions.
An unusual form of calcium carbonate is the hexahydrate ikaite, . Ikaite is stable only below 8 °C.
Preparation
The vast majority of calcium carbonate used in industry is extracted by mining or quarrying. Pure calcium carbonate (such as for food or pharmaceutical use), can be produced from a pure quarried source (usually marble).
Alternatively, calcium carbonate is prepared from calcium oxide. Water is added to give calcium hydroxide then carbon dioxide is passed through this solution to precipitate the desired calcium carbonate, referred to in the industry as precipitated calcium carbonate (PCC) This process is called carbonatation:
In a laboratory, calcium carbonate can easily be crystallized from calcium chloride (), by placing an aqueous solution of in a desiccator alongside ammonium carbonate . In the desiccator, ammonium carbonate is exposed to air and decomposes into ammonia, carbon dioxide, and water. The carbon dioxide then diffuses into the aqueous solution of calcium chloride, reacts with the calcium ions and the water, and forms calcium carbonate.
Structure
The thermodynamically stable form of under normal conditions is hexagonal β- (the mineral calcite). Other forms can be prepared, the denser (2.83 g/cm3) orthorhombic λ- (the mineral aragonite) and hexagonal μ-, occurring as the mineral vaterite. The aragonite form can be prepared by precipitation at temperatures above 85 °C; the vaterite form can be prepared by precipitation at 60 °C. Calcite contains calcium atoms coordinated by six oxygen atoms; in aragonite they are coordinated by nine oxygen atoms. The vaterite structure is not fully understood. Magnesium carbonate () has the calcite structure, whereas strontium carbonate () and barium carbonate () adopt the aragonite structure, reflecting their larger ionic radii.
Polymorphs
Calcium carbonate crystallizes in three anhydrous polymorphs, of which calcite is the thermodynamically most stable at room temperature, aragonite is only slightly less so, and vaterite is the least stable.
Crystal structure
The calcite crystal structure is trigonal, with space group Rc (No. 167 in the International Tables for Crystallography), and Pearson symbol hR10. Aragonite is orthorhombic, with space group Pmcn (No 62), and Pearson Symbol oP20. Vaterite is composed of at least two different coexisting crystallographic structures. The major structure exhibits hexagonal symmetry in space group P63/mmc, the minor structure is still unknown.
Crystallization
All three polymorphs crystallize simultaneously from aqueous solutions under ambient conditions. In additive-free aqueous solutions, calcite forms easily as the major product, while aragonite appears only as a minor product.
At high saturation, vaterite is typically the first phase precipitated, which is followed by a transformation of the vaterite to calcite. This behavior seems to follow Ostwald's rule, in which the least stable polymorph crystallizes first, followed by the crystallization of different polymorphs via a sequence of increasingly stable phases. However, aragonite, whose stability lies between those of vaterite and calcite, seems to be the exception to this rule, as aragonite does not form as a precursor to calcite under ambient conditions.
Aragonite occurs in majority when the reaction conditions inhibit the formation of calcite and/or promote the nucleation of aragonite. For example, the formation of aragonite is promoted by the presence of magnesium ions, or by using proteins and peptides derived from biological calcium carbonate. Some polyamines such as cadaverine and Poly(ethylene imine) have been shown to facilitate the formation of aragonite over calcite.
Selection by organisms
Organisms, such as molluscs and arthropods, have shown the ability to grow all three crystal polymorphs of calcium carbonate, mainly as protection (shells) and muscle attachments. Moreover, they exhibit a remarkable capability of phase selection over calcite and aragonite, and some organisms can switch between the two polymorphs. The ability of phase selection is usually attributed to the use of specific macromolecules or combinations of macromolecules by such organisms.
Occurrence
Geological sources
Calcite, aragonite and vaterite are pure calcium carbonate minerals. Industrially important source rocks which are predominantly calcium carbonate include limestone, chalk, marble and travertine.
Biological sources
Eggshells, snail shells and most seashells are predominantly calcium carbonate and can be used as industrial sources of that chemical. Oyster shells have enjoyed recent recognition as a source of dietary calcium, but are also a practical industrial source. Dark green vegetables such as broccoli and kale contain dietarily significant amounts of calcium carbonate, but they are not practical as an industrial source.
Annelids in the family Lumbricidae, earthworms, possess a regionalization of the digestive track called calciferous glands, Kalkdrüsen, or glandes de Morren, that processes calcium and into calcium carbonate, which is later excreted into the dirt. The function of these glands is unknown but is believed to serve as a regulation mechanism within the animals' tissues. This process is ecologically significant, stabilizing the pH of acid soils.
Extraterrestrial
Beyond Earth, strong evidence suggests the presence of calcium carbonate on Mars. Signs of calcium carbonate have been detected at more than one location (notably at Gusev and Huygens craters). This provides some evidence for the past presence of liquid water.
Geology
Carbonate is found frequently in geologic settings and constitutes an enormous carbon reservoir. Calcium carbonate occurs as aragonite, calcite and dolomite as significant constituents of the calcium cycle. The carbonate minerals form the rock types: limestone, chalk, marble, travertine, tufa, and others.
In warm, clear tropical waters corals are more abundant than towards the poles where the waters are cold. Calcium carbonate contributors, including plankton (such as coccoliths and planktic foraminifera), coralline algae, sponges, brachiopods, echinoderms, bryozoa and mollusks, are typically found in shallow water environments where sunlight and filterable food are more abundant. Cold-water carbonates do exist at higher latitudes but have a very slow growth rate. The calcification processes are changed by ocean acidification.
Where the oceanic crust is subducted under a continental plate sediments will be carried down to warmer zones in the asthenosphere and lithosphere. Under these conditions calcium carbonate decomposes to produce carbon dioxide which, along with other gases, give rise to explosive volcanic eruptions.
Carbonate compensation depth
The carbonate compensation depth (CCD) is the point in the ocean where the rate of precipitation of calcium carbonate is balanced by the rate of dissolution due to the conditions present. Deep in the ocean, the temperature drops and pressure increases. Increasing pressure also increases the solubility of calcium carbonate. Calcium carbonate is unusual in that its solubility increases with decreasing temperature. The carbonate compensation depth ranges from 4,000 to 6,000 meters below sea level in modern oceans, and the various polymorphs (calcite, aragonite) have different compensation depths based on their stability.
Role in taphonomy
Calcium carbonate can preserve fossils through permineralization. Most of the vertebrate fossils of the Two Medicine Formation—a geologic formation known for its duck-billed dinosaur eggs—are preserved by permineralization. This type of preservation conserves high levels of detail, even down to the microscopic level. However, it also leaves specimens vulnerable to weathering when exposed to the surface.
Trilobite populations were once thought to have composed the majority of aquatic life during the Cambrian, due to the fact that their calcium carbonate-rich shells were more easily preserved than those of other species, which had purely chitinous shells.
Uses
Construction
The main use of calcium carbonate is in the construction industry, either as a building material, or limestone aggregate for road building, as an ingredient of cement, or as the starting material for the preparation of builders' lime by burning in a kiln. However, because of weathering mainly caused by acid rain, calcium carbonate (in limestone form) is no longer used for building purposes on its own, but only as a raw primary substance for building materials.
Calcium carbonate is also used in the purification of iron from iron ore in a blast furnace. The carbonate is calcined in situ to give calcium oxide, which forms a slag with various impurities present, and separates from the purified iron.
In the oil industry, calcium carbonate is added to drilling fluids as a formation-bridging and filtercake-sealing agent; it is also a weighting material which increases the density of drilling fluids to control the downhole pressure. Calcium carbonate is added to swimming pools, as a pH corrector for maintaining alkalinity and offsetting the acidic properties of the disinfectant agent.
It is also used as a raw material in the refining of sugar from sugar beet; it is calcined in a kiln with anthracite to produce calcium oxide and carbon dioxide. This burnt lime is then slaked in fresh water to produce a calcium hydroxide suspension for the precipitation of impurities in raw juice during carbonatation.
Calcium carbonate in the form of chalk has traditionally been a major component of blackboard chalk. However, modern manufactured chalk is mostly gypsum, hydrated calcium sulfate . Calcium carbonate is a main source for growing biorock. Precipitated calcium carbonate (PCC), pre-dispersed in slurry form, is a common filler material for latex gloves with the aim of achieving maximum saving in material and production costs.
Fine ground calcium carbonate (GCC) is an essential ingredient in the microporous film used in diapers and some building films, as the pores are nucleated around the calcium carbonate particles during the manufacture of the film by biaxial stretching. GCC and PCC are used as a filler in paper because they are cheaper than wood fiber. Printing and writing paper can contain 10–20% calcium carbonate. In North America, calcium carbonate has begun to replace kaolin in the production of glossy paper. Europe has been practicing this as alkaline papermaking or acid-free papermaking for some decades. PCC used for paper filling and paper coatings is precipitated and prepared in a variety of shapes and sizes having characteristic narrow particle size distributions and equivalent spherical diameters of 0.4 to 3 micrometers.
Calcium carbonate is widely used as an extender in paints, in particular matte emulsion paint where typically 30% by weight of the paint is either chalk or marble. It is also a popular filler in plastics. Some typical examples include around 15–20% loading of chalk in unplasticized polyvinyl chloride (uPVC) drainpipes, 5–15% loading of stearate-coated chalk or marble in uPVC window profile. PVC cables can use calcium carbonate at loadings of up to 70 phr (parts per hundred parts of resin) to improve mechanical properties (tensile strength and elongation) and electrical properties (volume resistivity). Polypropylene compounds are often filled with calcium carbonate to increase rigidity, a requirement that becomes important at high usage temperatures. Here the percentage is often 20–40%. It also routinely used as a filler in thermosetting resins (sheet and bulk molding compounds) and has also been mixed with ABS, and other ingredients, to form some types of compression molded "clay" poker chips. Precipitated calcium carbonate, made by dropping calcium oxide into water, is used by itself or with additives as a white paint, known as whitewashing.
Calcium carbonate is added to a wide range of trade and do it yourself adhesives, sealants, and decorating fillers. Ceramic tile adhesives typically contain 70% to 80% limestone. Decorating crack fillers contain similar levels of marble or dolomite. It is also mixed with putty in setting stained glass windows, and as a resist to prevent glass from sticking to kiln shelves when firing glazes and paints at high temperature.
In ceramic glaze applications, calcium carbonate is known as whiting, and is a common ingredient for many glazes in its white powdered form. When a glaze containing this material is fired in a kiln, the whiting acts as a flux material in the glaze. Ground calcium carbonate is an abrasive (both as scouring powder and as an ingredient of household scouring creams), in particular in its calcite form, which has the relatively low hardness level of 3 on the Mohs scale, and will therefore not scratch glass and most other ceramics, enamel, bronze, iron, and steel, and have a moderate effect on softer metals like aluminium and copper. A paste made from calcium carbonate and deionized water can be used to clean tarnish on silver.
Health and diet
Calcium carbonate is widely used medicinally as an inexpensive dietary calcium supplement for gastric antacid (such as Tums and Eno). It may be used as a phosphate binder for the treatment of hyperphosphatemia (primarily in patients with chronic kidney failure). It is used in the pharmaceutical industry as an inert filler for tablets and other pharmaceuticals.
Calcium carbonate is used in the production of calcium oxide as well as toothpaste and has seen a resurgence as a food preservative and color retainer, when used in or with products such as organic apples.
Calcium carbonate is used therapeutically as phosphate binder in patients on maintenance haemodialysis. It is the most common form of phosphate binder prescribed, particularly in non-dialysis chronic kidney disease. Calcium carbonate is the most commonly used phosphate binder, but clinicians are increasingly prescribing the more expensive, non-calcium-based phosphate binders, particularly sevelamer.
Excess calcium from supplements, fortified food, and high-calcium diets can cause milk-alkali syndrome, which has serious toxicity and can be fatal. In 1915, Bertram Sippy introduced the "Sippy regimen" of hourly ingestion of milk and cream, and the gradual addition of eggs and cooked cereal, for 10 days, combined with alkaline powders, which provided symptomatic relief for peptic ulcer disease. Over the next several decades, the Sippy regimen resulted in kidney failure, alkalosis, and hypercalcaemia, mostly in men with peptic ulcer disease. These adverse effects were reversed when the regimen stopped, but it was fatal in some patients with protracted vomiting. Milk-alkali syndrome declined in men after effective treatments for peptic ulcer disease arose. Since the 1990s it has been most frequently reported in women taking calcium supplements above the recommended range of 1.2 to 1.5 grams daily, for prevention and treatment of osteoporosis, and is exacerbated by dehydration. Calcium has been added to over-the-counter products, which contributes to inadvertent excessive intake. Excessive calcium intake can lead to hypercalcemia, complications of which include vomiting, abdominal pain and altered mental status.
As a food additive it is designated E170, and it has an INS number of 170. Used as an acidity regulator, anticaking agent, stabilizer or color it is approved for usage in the EU, US and Australia and New Zealand. It is "added by law to all UK milled bread flour except wholemeal". It is used in some soy milk and almond milk products as a source of dietary calcium; at least one study suggests that calcium carbonate might be as bioavailable as the calcium in cow's milk. Calcium carbonate is also used as a firming agent in many canned and bottled vegetable products.
Several calcium supplement formulations have been documented to contain the chemical element lead, posing a public health concern. Lead is commonly found in natural sources of calcium.
Agriculture and aquaculture
Agricultural lime, powdered chalk or limestone, is used as a cheap method of neutralising acidic soil, making it suitable for planting, also used in aquaculture industry for pH regulation of pond soil before initiating culture. There is interest in understanding whether or not it can affect pesticide adsorption and desorption in calcareous soil.
Household cleaning
Calcium carbonate is a key ingredient in many household cleaning powders like Comet and is used as a scrubbing agent.
Pollution mitigation
In 1989, a researcher, Ken Simmons, introduced into the Whetstone Brook in Massachusetts. His hope was that the calcium carbonate would counter the acid in the stream from acid rain and save the trout that had ceased to spawn. Although his experiment was a success, it did increase the amount of aluminium ions in the area of the brook that was not treated with the limestone. This shows that can be added to neutralize the effects of acid rain in river ecosystems. Currently calcium carbonate is used to neutralize acidic conditions in both soil and water. Since the 1970s, such liming has been practiced on a large scale in Sweden to mitigate acidification and several thousand lakes and streams are limed repeatedly.
Calcium carbonate is also used in flue-gas desulfurization applications eliminating harmful and emissions from coal and other fossil fuels burnt in large fossil fuel power stations.
Plastics
Calcium carbonate is commonly used in the plastic industry as a filler. When it is incorporated in a plastic material, it can improve the hardness, stiffness, dimensional stability and processability of the material.
Calcination equilibrium
Calcination of limestone using charcoal fires to produce quicklime has been practiced since antiquity by cultures all over the world. The temperature at which limestone yields calcium oxide is usually given as 825 °C, but stating an absolute threshold is misleading. Calcium carbonate exists in equilibrium with calcium oxide and carbon dioxide at any temperature. At each temperature there is a partial pressure of carbon dioxide that is in equilibrium with calcium carbonate. At room temperature the equilibrium overwhelmingly favors calcium carbonate, because the equilibrium pressure is only a tiny fraction of the partial pressure in air, which is about 0.035 kPa.
At temperatures above 550 °C the equilibrium pressure begins to exceed the pressure in air. So above 550 °C, calcium carbonate begins to outgas into air. However, in a charcoal fired kiln, the concentration of will be much higher than it is in air. Indeed, if all the oxygen in the kiln is consumed in the fire, then the partial pressure of in the kiln can be as high as 20 kPa.
The table shows that this partial pressure is not achieved until the temperature is nearly 800 °C. For the outgassing of from calcium carbonate to happen at an economically useful rate, the equilibrium pressure must significantly exceed the ambient pressure of . And for it to happen rapidly, the equilibrium pressure must exceed total atmospheric pressure of 101 kPa, which happens at 898 °C.
{| class="wikitable"
|+ Equilibrium pressure of over (P) versus temperature (T).
|-
!P (kPa)
|0.055||0.13||0.31||1.80||5.9||9.3||14||24||34||51||72 ||80||91||101||179||901||3961
|-
!T (°C)
|550||587||605||680||727||748||777||800||830||852||871||881||891||898||937||1082||1241
|}
Solubility
With varying pressure
Calcium carbonate is poorly soluble in pure water (47 mg/L at normal atmospheric partial pressure as shown below).
The equilibrium of its solution is given by the equation (with dissolved calcium carbonate on the right):
{| width="500"
| style="width:50%; height:30px;"|
| Ksp = to at 25 °C
|}
where the solubility product for is given as anywhere from Ksp = to Ksp = at 25 °C, depending upon the data source. What the equation means is that the product of molar concentration of calcium ions (moles of dissolved per liter of solution) with the molar concentration of dissolved cannot exceed the value of Ksp. This seemingly simple solubility equation, however, must be taken along with the more complicated equilibrium of carbon dioxide with water (see carbonic acid). Some of the combines with in the solution according to
{| width="500"
| style="width:50%; height:25px;"|
| Ka2 = at 25 °C
|}
is known as the bicarbonate ion. Calcium bicarbonate is many times more soluble in water than calcium carbonate—indeed it exists only in solution.
Some of the combines with in solution according to
{|
| width=250|
| Ka1 = at 25 °C
|}
Some of the breaks up into water and dissolved carbon dioxide according to
{|
|width=250|
| Kh = at 25 °C
|}
And dissolved carbon dioxide is in equilibrium with atmospheric carbon dioxide according to
{|
| width=250| =
| where = 29.76 atm/(mol/L) at 25 °C (Henry volatility), and P is the partial pressure.
|}
For ambient air, P is around atm (or equivalently 35 Pa). The last equation above fixes the concentration of dissolved as a function of P, independent of the concentration of dissolved . At atmospheric partial pressure of , dissolved concentration is moles per liter. The equation before that fixes the concentration of as a function of concentration. For [] = , it results in = moles per liter. When is known, the remaining three equations together with
{| width="450"
| style="width:50%; height:25px;"|
| K = 10−14 at 25 °C
|}
(which is true for all aqueous solutions), and the constraint that the solution must be electrically neutral, i.e., the overall charge of dissolved positive ions must be cancelled out by the overall charge of dissolved negative ions , make it possible to solve simultaneously for the remaining five unknown concentrations (the previously mentioned form of the neutrality is valid only if calcium carbonate has been put in contact with pure water or with a neutral pH solution; in the case where the initial water solvent pH is not neutral, the balance is not neutral).
The adjacent table shows the result for and (in the form of pH) as a function of ambient partial pressure of (Ksp = has been taken for the calculation).
At atmospheric levels of ambient the table indicates that the solution will be slightly alkaline with a maximum solubility of 47 mg/L.
As ambient partial pressure is reduced below atmospheric levels, the solution becomes more and more alkaline. At extremely low P, dissolved , bicarbonate ion, and carbonate ion largely evaporate from the solution, leaving a highly alkaline solution of calcium hydroxide, which is more soluble than . For P = 10−12 atm, the product is still below the solubility product of (). For still lower pressure, precipitation will occur before precipitation.
As ambient partial pressure increases to levels above atmospheric, pH drops, and much of the carbonate ion is converted to bicarbonate ion, which results in higher solubility of .
The effect of the latter is especially evident in day-to-day life of people who have hard water. Water in aquifers underground can be exposed to levels of much higher than atmospheric. As such water percolates through calcium carbonate rock, the dissolves according to one of the trends above. When that same water then emerges from the tap, in time it comes into equilibrium with levels in the air by outgassing its excess . The calcium carbonate becomes less soluble as a result, and the excess precipitates as lime scale. This same process is responsible for the formation of stalactites and stalagmites in limestone caves.
Two hydrated phases of calcium carbonate, monohydrocalcite and ikaite , may precipitate from water at ambient conditions and persist as metastable phases.
With varying pH, temperature and salinity: scaling in swimming pools
In contrast to the open equilibrium scenario above, many swimming pools are managed by addition of sodium bicarbonate () to the concentration of about 2 mmol/L as a buffer, then control of pH through use of HCl, , , NaOH or chlorine formulations that are acidic or basic. In this situation, dissolved inorganic carbon (total inorganic carbon) is far from equilibrium with atmospheric . Progress towards equilibrium through outgassing of is slowed by
In this situation, the dissociation constants for the much faster reactions
allow the prediction of concentrations of each dissolved inorganic carbon species in solution, from the added concentration of (which constitutes more than 90% of Bjerrum plot species from pH 7 to pH 8 at 25 °C in fresh water). Addition of will increase concentration at any pH. Rearranging the equations given above, we can see that = , and [] = . Therefore, when concentration is known, the maximum concentration of ions before scaling through precipitation can be predicted from the formula:
[]max = ×
The solubility product for (Ksp) and the dissociation constants for the dissolved inorganic carbon species (including Ka2) are all substantially affected by temperature and salinity, with the overall effect that []max increases from freshwater to saltwater, and decreases with rising temperature, pH, or added bicarbonate level, as illustrated in the accompanying graphs.
The trends are illustrative for pool management, but whether scaling occurs also depends on other factors including interactions with , and other ions in the pool, as well as supersaturation effects. Scaling is commonly observed in electrolytic chlorine generators, where there is a high pH near the cathode surface and scale deposition further increases temperature. This is one reason that some pool operators prefer borate over bicarbonate as the primary pH buffer, and avoid the use of pool chemicals containing calcium.
Solubility in a strong or weak acid solution
Solutions of strong (HCl), moderately strong (sulfamic) or weak (acetic, citric, sorbic, lactic, phosphoric) acids are commercially available. They are commonly used as descaling agents to remove limescale deposits. The maximum amount of that can be "dissolved" by one liter of an acid solution can be calculated using the above equilibrium equations.
In the case of a strong monoacid with decreasing acid concentration [A] = [], we obtain (with molar mass = 100 g/mol):
{| class="wikitable"
! [A] (mol/L)
| 1
| 10−1
| 10−2
| 10−3
| 10−4
| 10−5
| 10−6
| 10−7
| 10−10
|-
! Initial pH
| 0.00||1.00||2.00||3.00||4.00||5.00||6.00||6.79||7.00
|-
!Final pH
| 6.75||7.25||7.75||8.14||8.25||8.26||8.26||8.26||8.27
|-
! Dissolved (g/L of acid)
| 50.0||5.00||0.514||0.0849||0.0504||0.0474||0.0471||0.0470||0.0470
|}
where the initial state is the acid solution with no (not taking into account possible dissolution) and the final state is the solution with saturated . For strong acid concentrations, all species have a negligible concentration in the final state with respect to and so that the neutrality equation reduces approximately to 2[] = [] yielding [] ≈ 0.5 []. When the concentration decreases, [] becomes non-negligible so that the preceding expression is no longer valid. For vanishing acid concentrations, one can recover the final pH and the solubility of in pure water.
In the case of a weak monoacid (here we take acetic acid with pKa = 4.76) with decreasing total acid concentration [A] = [] + [AH], we obtain:
{| class="wikitable"
![A] (mol/L)
| [] ≈ 0.5 []
| 10−1
| 10−2
| 10−3
| 10−4
| 10−5
| 10−6
| 10−7
| 10−10
|-
!Initial pH
| 2.38||2.88||3.39||3.91||4.47||5.15||6.02||6.79||7.00
|-
!Final pH
| 6.75||7.25||7.75||8.14||8.25||8.26||8.26||8.26||8.27
|-
!|Dissolved (g/L of acid)
| 49.5||4.99||0.513||0.0848||0.0504||0.0474||0.0471||0.0470||0.0470
|}
For the same total acid concentration, the initial pH of the weak acid is less acid than the one of the strong acid; however, the maximum amount of which can be dissolved is approximately the same. This is because in the final state, the pH is larger than the pKa, so that the weak acid is almost completely dissociated, yielding in the end as many ions as the strong acid to "dissolve" the calcium carbonate.
The calculation in the case of phosphoric acid (which is the most widely used for domestic applications) is more complicated since the concentrations of the four dissociation states corresponding to this acid must be calculated together with [], [], [], [] and []. The system may be reduced to a seventh degree equation for [] the numerical solution of which gives
{| class="wikitable"
! [A] (mol/L)
| 1
| 10−1
| 10−2
| 10−3
| 10−4
| 10−5
| 10−6
| 10−7
| 10−10
|-
! Initial pH
| 1.08||1.62||2.25||3.05||4.01||5.00||5.97||6.74||7.00
|-
! Final pH
| 6.71||7.17||7.63||8.06||8.24||8.26||8.26||8.26||8.27
|-
! Dissolved (g/L of acid)
| 62.0||7.39||0.874||0.123||0.0536||0.0477||0.0471||0.0471||0.0470
|}
where [A] = is the total acid concentration. Thus phosphoric acid is more efficient than a monoacid since at the final almost neutral pH, the second dissociated state concentration [] is not negligible (see phosphoric acid).
| Physical sciences | Salts | null |
44734 | https://en.wikipedia.org/wiki/Chalk | Chalk | Chalk is a soft, white, porous, sedimentary carbonate rock. It is a form of limestone composed of the mineral calcite and originally formed deep under the sea by the compression of microscopic plankton that had settled to the sea floor. Chalk is common throughout Western Europe, where deposits underlie parts of France, and steep cliffs are often seen where they meet the sea in places such as the Dover cliffs on the Kent coast of the English Channel.
Chalk is mined for use in industry, such as for quicklime, bricks and builder's putty, and in agriculture, for raising pH in soils with high acidity. It is also used for "blackboard chalk" for writing and drawing on various types of surfaces, although these can also be manufactured from other carbonate-based minerals, or gypsum.
Description
Chalk is a fine-textured, earthy type of limestone distinguished by its light colour, softness, and high porosity. It is composed mostly of tiny fragments of the calcite shells or skeletons of plankton, such as foraminifera or coccolithophores. These fragments mostly take the form of calcite plates ranging from 0.5 to 4 microns in size, though about 10% to 25% of a typical chalk is composed of fragments that are 10 to 100 microns in size. The larger fragments include intact plankton skeletons and skeletal fragments of larger organisms, such as molluscs, echinoderms, or bryozoans.
Chalk is typically almost pure calcite, , with just 2% to 4% of other minerals. These are usually quartz and clay minerals, though collophane (cryptocrystalline apatite, a phosphate mineral) is also sometimes present, as nodules or as small pellets interpreted as fecal pellets. In some chalk beds, the calcite has been converted to dolomite, , and in a few cases the dolomitized chalk has been dedolomitized back to calcite.
Chalk is highly porous, with typical values of porosity ranging from 35 to 47 per cent. While it is similar in appearance to both gypsum and diatomite, chalk is identifiable by its hardness, fossil content, and its reaction to acid (it produces effervescence on contact).
Formation
In Western Europe, chalk was formed in the Late Cretaceous Epoch and the early Palaeocene Epoch (between 100 and 61 million years ago). It was deposited on extensive continental shelves at depths between , during a time of nonseasonal (likely arid) climate that reduced the amount of erosion from nearby exposed rock. The lack of nearby erosion explains the high purity of chalk. The coccolithophores, foraminifera, and other microscopic organisms from which the chalk came mostly form low-magnesium calcite skeletons, so the sediments were already in the form of highly stable low-magnesium calcite when deposited. This is in contrast with most other limestones, which formed from high-magnesium calcite or aragonite that rapidly converted to the more stable low-magnesium calcite after deposition, resulting in the early cementation of such limestones. In chalk, the absence of calcium carbonate conversion process prevented early cementation, and it accounts for chalk's high porosity. Additionally, chalk is the only form of limestone that commonly shows signs of compaction.
Flint (a type of chert) is very common as bands parallel to the bedding or as nodules in seams, or linings to fractures, embedded in chalk. It is probably derived from sponge spicules or other siliceous organisms as water is expelled upwards during compaction. Flint is often deposited around larger fossils such as Echinoidea which may be silicified (i.e. replaced molecule by molecule by flint).
Geology and geographic distribution
Chalk is so common in Cretaceous marine beds that the Cretaceous Period was named for these deposits. The name Cretaceous was derived from Latin creta, meaning chalk. Some deposits of chalk were formed after the Cretaceous.
The Chalk Group is a European stratigraphic unit deposited during the late Cretaceous Period. It forms the famous White Cliffs of Dover in Kent, England, as well as their counterparts of the Cap Blanc Nez on the other side of the Dover Strait. The Champagne region of France is mostly underlain by chalk deposits, which contain artificial caves used for wine storage. Some of the highest chalk cliffs in the world occur at Jasmund National Park in Germany and at Møns Klint in Denmark.
Chalk deposits are also found in Cretaceous beds on other continents, such as the Austin Chalk, Selma Group, and Niobrara Formations of the North American interior. Chalk is also found in western Egypt (Khoman Formation) and western Australia (Miria Formation).
Chalk of Oligocene to Neogene age has been found in drill cores of rock under the Pacific Ocean at Stewart Arch in the Solomon Islands.
There are layers of chalk, containing Globorotalia, in the Nicosia Formation of Cyprus, which formed during the Pliocene.
Mining
Chalk is mined from chalk deposits both above ground and underground. Chalk mining boomed during the Industrial Revolution, due to the need for chalk products such as quicklime and bricks.
Uses
Most people first encounter chalk in school where it refers to blackboard chalk, which was originally made of mineral chalk, since it readily crumbles and leaves particles that stick loosely to rough surfaces, allowing it to make writing that can be readily erased. Blackboard chalk manufacturers now may use mineral chalk, other mineral sources of calcium carbonate, or the mineral gypsum (calcium sulfate). While gypsum-based blackboard chalk is the lowest cost to produce, and thus widely used in the developing world, use of carbonate-based chalk produces larger particles and thus less dust, and it is marketed as "dustless chalk".
Coloured chalks, pastel chalks, and sidewalk chalk (shaped into larger sticks and often coloured), used to draw on sidewalks, streets, and driveways, are primarily made of gypsum rather than calcium carbonate chalk.
Magnesium carbonate chalk is commonly used as a drying agent to obtain better grip by gymnasts and rock climbers.
Glazing putty mainly contains chalk as a filler in linseed oil.
Chalk and other forms of limestone may be used for their properties as a base. Chalk is a source of quicklime by thermal decomposition, or slaked lime following quenching of quicklime with water. In agriculture, chalk is used for raising pH in soils with high acidity. Small doses of chalk can also be used as an antacid. Additionally, the small particles of chalk make it a substance ideal for cleaning and polishing. For example, toothpaste commonly contains small amounts of chalk, which serves as a mild abrasive. Polishing chalk is chalk prepared with a carefully controlled grain size, for very fine polishing of metals.
French chalk (also known as tailor's chalk) is traditionally a hard chalk used to make temporary markings on cloth, mainly by tailors. It is now usually made of talc (magnesium silicate).
Chalk beds form important petroleum reservoirs in the North Sea and along the Gulf Coast of North America.
Previous uses
In southeast England, deneholes are a notable example of ancient chalk pits. Such bell pits may also mark the sites of ancient flint mines, where the prime object was to remove flint nodules for stone tool manufacture. The surface remains at Cissbury are one such example, but perhaps the most famous is the extensive complex at Grimes Graves in Norfolk.
Chalk was traditionally used in recreation. In field sports, such as tennis played on grass, powdered chalk was used to mark the boundary lines of the playing field or court. If a ball hits the line, a cloud of chalk or pigment dust will be visible. In recent years, powdered chalk has been replaced with titanium dioxide. In gymnastics, rock-climbing, weightlifting and tug of war, chalk — now usually magnesium carbonate — is applied to the hands and feet to remove perspiration and reduce slipping.
Chalk may also be used as a house construction material instead of brick or wattle and daub: quarried chalk was cut into blocks and used as ashlar, or loose chalk was rammed into blocks and laid in mortar. There are still houses standing which have been constructed using chalk as the main building material. Most are pre-Victorian though a few are more recent.
A mixture of chalk and mercury can be used as fingerprint powder. However, because of the toxicity of the mercury, the use of such mixtures for fingerprinting was abandoned in 1967.
| Physical sciences | Sedimentary rocks | Earth science |
44758 | https://en.wikipedia.org/wiki/Goldbach%27s%20conjecture | Goldbach's conjecture | Goldbach's conjecture is one of the oldest and best-known unsolved problems in number theory and all of mathematics. It states that every even natural number greater than 2 is the sum of two prime numbers.
The conjecture has been shown to hold for all integers less than but remains unproven despite considerable effort.
History
Origins
On 7 June 1742, the Prussian mathematician Christian Goldbach wrote a letter to Leonhard Euler (letter XLIII), in which he proposed the following conjecture:
Goldbach was following the now-abandoned convention of considering 1 to be a prime number, so that a sum of units would be a sum of primes.
He then proposed a second conjecture in the margin of his letter, which implies the first:
Euler replied in a letter dated 30 June 1742 and reminded Goldbach of an earlier conversation they had had (""), in which Goldbach had remarked that the first of those two conjectures would follow from the statement
This is in fact equivalent to his second, marginal conjecture.
In the letter dated 30 June 1742, Euler stated:
Similar conjecture by Descartes
René Descartes wrote that "Every even number can be expressed as the sum of at most three primes." The proposition is equivalent to Goldbach's conjecture, and Paul Erdős said that "Descartes actually discovered this before Goldbach... but it is better that the conjecture was named for Goldbach because, mathematically speaking, Descartes was infinitely rich and Goldbach was very poor."
Partial results
The strong Goldbach conjecture is much more difficult than the weak Goldbach conjecture, which says that every odd integer greater than 5 is the sum of three primes. Using Vinogradov's method, Nikolai Chudakov, Johannes van der Corput, and Theodor Estermann showed (1937–1938) that almost all even numbers can be written as the sum of two primes (in the sense that the fraction of even numbers up to some which can be so written tends towards 1 as increases). In 1930, Lev Schnirelmann proved that any natural number greater than 1 can be written as the sum of not more than prime numbers, where is an effectively computable constant; see Schnirelmann density. Schnirelmann's constant is the lowest number with this property. Schnirelmann himself obtained . This result was subsequently enhanced by many authors, such as Olivier Ramaré, who in 1995 showed that every even number is in fact the sum of at most 6 primes. The best known result currently stems from the proof of the weak Goldbach conjecture by Harald Helfgott, which directly implies that every even number is the sum of at most 4 primes.
In 1924, Hardy and Littlewood showed under the assumption of the generalized Riemann hypothesis that the number of even numbers up to violating the Goldbach conjecture is much less than for small .
In 1948, using sieve theory methods, Alfréd Rényi showed that every sufficiently large even number can be written as the sum of a prime and an almost prime with at most factors. Chen Jingrun showed in 1973 using sieve theory that every sufficiently large even number can be written as the sum of either two primes, or a prime and a semiprime (the product of two primes). See Chen's theorem for further information.
In 1975, Hugh Lowell Montgomery and Bob Vaughan showed that "most" even numbers are expressible as the sum of two primes. More precisely, they showed that there exist positive constants and such that for all sufficiently large numbers , every even number less than is the sum of two primes, with at most exceptions. In particular, the set of even integers that are not the sum of two primes has density zero.
In 1951, Yuri Linnik proved the existence of a constant such that every sufficiently large even number is the sum of two primes and at most powers of 2. János Pintz and Imre Ruzsa found in 2020 that works. Assuming the generalized Riemann hypothesis, also works, as shown by Roger Heath-Brown and Jan-Christoph Schlage-Puchta in 2002.
A proof for the weak conjecture was submitted in 2013 by Harald Helfgott to Annals of Mathematics Studies series. Although the article was accepted, Helfgott decided to undertake the major modifications suggested by the referee. Despite several revisions, Helfgott's proof has not yet appeared in a peer-reviewed publication. The weak conjecture is implied by the strong conjecture, as if is a sum of two primes, then is a sum of three primes. However, the converse implication and thus the strong Goldbach conjecture would remain unproven if Helfgott's proof is correct.
Computational results
For small values of , the strong Goldbach conjecture (and hence the weak Goldbach conjecture) can be verified directly. For instance, in 1938, Nils Pipping laboriously verified the conjecture up to . With the advent of computers, many more values of have been checked; T. Oliveira e Silva ran a distributed computer search that has verified the conjecture for (and double-checked up to ) as of 2013. One record from this search is that is the smallest number that cannot be written as a sum of two primes where one is smaller than 9781.
In popular culture
Goldbach's Conjecture () is the title of the biography of Chinese mathematician and number theorist Chen Jingrun, written by Xu Chi.
The conjecture is a central point in the plot of the 1992 novel Uncle Petros and Goldbach's Conjecture by Greek author Apostolos Doxiadis, in the short story "Sixty Million Trillion Combinations" by Isaac Asimov and also in the 2008 mystery novel No One You Know by Michelle Richmond.
Goldbach's conjecture is part of the plot of the 2007 Spanish film Fermat's Room.
Goldbach's conjecture is featured as the main topic of research of actress Ella Rumpf's character Marguerite in the 2023 French-Swiss film Marguerite's Theorem.
Formal statement
Each of the three conjectures has a natural analog in terms of the modern definition of a prime, under which 1 is excluded. A modern version of the first conjecture is:
A modern version of the marginal conjecture is:
And a modern version of Goldbach's older conjecture of which Euler reminded him is:
These modern versions might not be entirely equivalent to the corresponding original statements. For example, if there were an even integer larger than 4, for a prime, that could not be expressed as the sum of two primes in the modern sense, then it would be a counterexample to the modern version of the third conjecture (without being a counterexample to the original version). The modern version is thus probably stronger (but in order to confirm that, one would have to prove that the first version, freely applied to any positive even integer , could not possibly rule out the existence of such a specific counterexample ). In any case, the modern statements have the same relationships with each other as the older statements did. That is, the second and third modern statements are equivalent, and either implies the first modern statement.
The third modern statement (equivalent to the second) is the form in which the conjecture is usually expressed today. It is also known as the "strong", "even", or "binary" Goldbach conjecture. A weaker form of the second modern statement, known as "Goldbach's weak conjecture", the "odd Goldbach conjecture", or the "ternary Goldbach conjecture", asserts that
Heuristic justification
Statistical considerations that focus on the probabilistic distribution of prime numbers present informal evidence in favour of the conjecture (in both the weak and strong forms) for sufficiently large integers: the greater the integer, the more ways there are available for that number to be represented as the sum of two or three other numbers, and the more "likely" it becomes that at least one of these representations consists entirely of primes.
A very crude version of the heuristic probabilistic argument (for the strong form of the Goldbach conjecture) is as follows. The prime number theorem asserts that an integer selected at random has roughly a chance of being prime. Thus if is a large even integer and is a number between 3 and , then one might expect the probability of and simultaneously being prime to be . If one pursues this heuristic, one might expect the total number of ways to write a large even integer as the sum of two odd primes to be roughly
Since , this quantity goes to infinity as increases, and one would expect that every large even integer has not just one representation as the sum of two primes, but in fact very many such representations.
This heuristic argument is actually somewhat inaccurate because it assumes that the events of and being prime are statistically independent of each other. For instance, if is odd, then is also odd, and if is even, then is even, a non-trivial relation because, besides the number 2, only odd numbers can be prime. Similarly, if is divisible by 3, and was already a prime other than 3, then would also be coprime to 3 and thus be slightly more likely to be prime than a general number. Pursuing this type of analysis more carefully, G. H. Hardy and John Edensor Littlewood in 1923 conjectured (as part of their Hardy–Littlewood prime tuple conjecture) that for any fixed , the number of representations of a large integer as the sum of primes with should be asymptotically equal to
where the product is over all primes , and is the number of solutions to the equation in modular arithmetic, subject to the constraints . This formula has been rigorously proven to be asymptotically valid for from the work of Ivan Matveevich Vinogradov, but is still only a conjecture when . In the latter case, the above formula simplifies to 0 when is odd, and to
when is even, where is Hardy–Littlewood's twin prime constant
This is sometimes known as the extended Goldbach conjecture. The strong Goldbach conjecture is in fact very similar to the twin prime conjecture, and the two conjectures are believed to be of roughly comparable difficulty.
The Goldbach partition function is the function that associates to each even integer the number of ways it can be decomposed into a sum of two primes. Its graph looks like a comet and is therefore called Goldbach's comet.
Goldbach's comet suggests tight upper and lower bounds on the number of representations of an even number as the sum of two primes, and also that the number of these representations depend strongly on the value modulo 3 of the number.
Related problems
Although Goldbach's conjecture implies that every positive integer greater than one can be written as a sum of at most three primes, it is not always possible to find such a sum using a greedy algorithm that uses the largest possible prime at each step. The Pillai sequence tracks the numbers requiring the largest number of primes in their greedy representations.
Similar problems to Goldbach's conjecture exist in which primes are replaced by other particular sets of numbers, such as the squares:
It was proven by Lagrange that every positive integer is the sum of four squares. See Waring's problem and the related Waring–Goldbach problem on sums of powers of primes.
Hardy and Littlewood listed as their Conjecture I: "Every large odd number () is the sum of a prime and the double of a prime". This conjecture is known as Lemoine's conjecture and is also called Levy's conjecture.
The Goldbach conjecture for practical numbers, a prime-like sequence of integers, was stated by Margenstern in 1984, and proved by Melfi in 1996: every even number is a sum of two practical numbers.
Harvey Dubner proposed a strengthening of the Goldbach conjecture that states that every even integer greater than 4208 is the sum of two twin primes (not necessarily belonging to the same pair). Only 34 even integers less than 4208 are not the sum of two twin primes; Dubner has verified computationally that this list is complete up to A proof of this stronger conjecture would not only imply Goldbach's conjecture, but also the twin prime conjecture.
According to Bertrand's postulate, for every integer , there is always at least one prime such that If the postulate were false, there would exist some integer for which no prime numbers lie between and , making it impossible to express as a sum of two primes.
Goldbach's conjecture is used when studying computation complexity. The connection is made through the Busy Beaver function, where BB(n) is the maximum number of steps taken by any n state Turing machine that halts. There is a 27-state Turing machine that halts if and only if Goldbach's conjecture is false. Hence if BB(27) was known, and the Turing machine did not stop in that number of steps, it would be known to run forever and hence no counterexamples exist (which proves the conjecture true). This is a completely impractical way to settle the conjecture; instead it is used to suggest that BB(27) will be very hard to compute, at least as difficult as settling the Goldbach conjecture.
| Mathematics | Sums and products | null |
44790 | https://en.wikipedia.org/wiki/Luminosity | Luminosity | Luminosity is an absolute measure of radiated electromagnetic energy per unit time, and is synonymous with the radiant power emitted by a light-emitting object. In astronomy, luminosity is the total amount of electromagnetic energy emitted per unit of time by a star, galaxy, or other astronomical objects.
In SI units, luminosity is measured in joules per second, or watts. In astronomy, values for luminosity are often given in the terms of the luminosity of the Sun, L⊙. Luminosity can also be given in terms of the astronomical magnitude system: the absolute bolometric magnitude (Mbol) of an object is a logarithmic measure of its total energy emission rate, while absolute magnitude is a logarithmic measure of the luminosity within some specific wavelength range or filter band.
In contrast, the term brightness in astronomy is generally used to refer to an object's apparent brightness: that is, how bright an object appears to an observer. Apparent brightness depends on both the luminosity of the object and the distance between the object and observer, and also on any absorption of light along the path from object to observer. Apparent magnitude is a logarithmic measure of apparent brightness. The distance determined by luminosity measures can be somewhat ambiguous, and is thus sometimes called the luminosity distance.
Measurement
When not qualified, the term "luminosity" means bolometric luminosity, which is measured either in the SI units, watts, or in terms of solar luminosities (). A bolometer is the instrument used to measure radiant energy over a wide band by absorption and measurement of heating. A star also radiates neutrinos, which carry off some energy (about 2% in the case of the Sun), contributing to the star's total luminosity. The IAU has defined a nominal solar luminosity of to promote publication of consistent and comparable values in units of the solar luminosity.
While bolometers do exist, they cannot be used to measure even the apparent brightness of a star because they are insufficiently sensitive across the electromagnetic spectrum and because most wavelengths do not reach the surface of the Earth. In practice bolometric magnitudes are measured by taking measurements at certain wavelengths and constructing a model of the total spectrum that is most likely to match those measurements. In some cases, the process of estimation is extreme, with luminosities being calculated when less than 1% of the energy output is observed, for example with a hot Wolf-Rayet star observed only in the infrared. Bolometric luminosities can also be calculated using a bolometric correction to a luminosity in a particular passband.
The term luminosity is also used in relation to particular passbands such as a visual luminosity of K-band luminosity. These are not generally luminosities in the strict sense of an absolute measure of radiated power, but absolute magnitudes defined for a given filter in a photometric system. Several different photometric systems exist. Some such as the UBV or Johnson system are defined against photometric standard stars, while others such as the AB system are defined in terms of a spectral flux density.
Stellar luminosity
A star's luminosity can be determined from two stellar characteristics: size and effective temperature. The former is typically represented in terms of solar radii, R⊙, while the latter is represented in kelvins, but in most cases neither can be measured directly. To determine a star's radius, two other metrics are needed: the star's angular diameter and its distance from Earth. Both can be measured with great accuracy in certain cases, with cool supergiants often having large angular diameters, and some cool evolved stars having masers in their atmospheres that can be used to measure the parallax using VLBI. However, for most stars the angular diameter or parallax, or both, are far below our ability to measure with any certainty. Since the effective temperature is merely a number that represents the temperature of a black body that would reproduce the luminosity, it obviously cannot be measured directly, but it can be estimated from the spectrum.
An alternative way to measure stellar luminosity is to measure the star's apparent brightness and distance. A third component needed to derive the luminosity is the degree of interstellar extinction that is present, a condition that usually arises because of gas and dust present in the interstellar medium (ISM), the Earth's atmosphere, and circumstellar matter. Consequently, one of astronomy's central challenges in determining a star's luminosity is to derive accurate measurements for each of these components, without which an accurate luminosity figure remains elusive. Extinction can only be measured directly if the actual and observed luminosities are both known, but it can be estimated from the observed colour of a star, using models of the expected level of reddening from the interstellar medium.
In the current system of stellar classification, stars are grouped according to temperature, with the massive, very young and energetic Class O stars boasting temperatures in excess of 30,000 K while the less massive, typically older Class M stars exhibit temperatures less than 3,500 K. Because luminosity is proportional to temperature to the fourth power, the large variation in stellar temperatures produces an even vaster variation in stellar luminosity. Because the luminosity depends on a high power of the stellar mass, high mass luminous stars have much shorter lifetimes. The most luminous stars are always young stars, no more than a few million years for the most extreme. In the Hertzsprung–Russell diagram, the x-axis represents temperature or spectral type while the y-axis represents luminosity or magnitude. The vast majority of stars are found along the main sequence with blue Class O stars found at the top left of the chart while red Class M stars fall to the bottom right. Certain stars like Deneb and Betelgeuse are found above and to the right of the main sequence, more luminous or cooler than their equivalents on the main sequence. Increased luminosity at the same temperature, or alternatively cooler temperature at the same luminosity, indicates that these stars are larger than those on the main sequence and they are called giants or supergiants.
Blue and white supergiants are high luminosity stars somewhat cooler than the most luminous main sequence stars. A star like Deneb, for example, has a luminosity around 200,000 L⊙, a spectral type of A2, and an effective temperature around 8,500 K, meaning it has a radius around . For comparison, the red supergiant Betelgeuse has a luminosity around 100,000 L⊙, a spectral type of M2, and a temperature around 3,500 K, meaning its radius is about . Red supergiants are the largest type of star, but the most luminous are much smaller and hotter, with temperatures up to 50,000 K and more and luminosities of several million L⊙, meaning their radii are just a few tens of R⊙. For example, R136a1 has a temperature over 46,000 K and a luminosity of more than 6,100,000 L⊙ (mostly in the UV), it is only .
Radio luminosity
The luminosity of a radio source is measured in , to avoid having to specify a bandwidth over which it is measured. The observed strength, or flux density, of a radio source is measured in Jansky where .
For example, consider a 10W transmitter at a distance of 1 million metres, radiating over a bandwidth of 1 MHz. By the time that power has reached the observer, the power is spread over the surface of a sphere with area or about , so its flux density is .
More generally, for sources at cosmological distances, a k-correction must be made for the spectral index α of the source, and a relativistic correction must be made for the fact that the frequency scale in the emitted rest frame is different from that in the observer's rest frame. So the full expression for radio luminosity, assuming isotropic emission, is
where Lν is the luminosity in , Sobs is the observed flux density in , DL is the luminosity distance in metres, z is the redshift, α is the spectral index (in the sense , and in radio astronomy, assuming thermal emission the spectral index is typically equal to 2.)
For example, consider a 1 Jy signal from a radio source at a redshift of 1, at a frequency of 1.4 GHz.
Ned Wright's cosmology calculator calculates a luminosity distance for a redshift of 1 to be 6701 Mpc = 2×1026 m giving a radio luminosity of .
To calculate the total radio power, this luminosity must be integrated over the bandwidth of the emission. A common assumption is to set the bandwidth to the observing frequency, which effectively assumes the power radiated has uniform intensity from zero frequency up to the observing frequency. In the case above, the total power is . This is sometimes expressed in terms of the total (i.e. integrated over all wavelengths) luminosity of the Sun which is , giving a radio power of .
Luminosity formulae
The Stefan–Boltzmann equation applied to a black body gives the value for luminosity for a black body, an idealized object which is perfectly opaque and non-reflecting:
where A is the surface area, T is the temperature (in kelvins) and is the Stefan–Boltzmann constant, with a value of
Imagine a point source of light of luminosity that radiates equally in all directions. A hollow sphere centered on the point would have its entire interior surface illuminated. As the radius increases, the surface area will also increase, and the constant luminosity has more surface area to illuminate, leading to a decrease in observed brightness.
where
is the area of the illuminated surface.
is the flux density of the illuminated surface.
The surface area of a sphere with radius r is , so for stars and other point sources of light:
where is the distance from the observer to the light source.
For stars on the main sequence, luminosity is also related to mass approximately as below:
Relationship to magnitude
Luminosity is an intrinsic measurable property of a star independent of distance. The concept of magnitude, on the other hand, incorporates distance. The apparent magnitude is a measure of the diminishing flux of light as a result of distance according to the inverse-square law. The Pogson logarithmic scale is used to measure both apparent and absolute magnitudes, the latter corresponding to the brightness of a star or other celestial body as seen if it would be located at an interstellar distance of . In addition to this brightness decrease from increased distance, there is an extra decrease of brightness due to extinction from intervening interstellar dust.
By measuring the width of certain absorption lines in the stellar spectrum, it is often possible to assign a certain luminosity class to a star without knowing its distance. Thus a fair measure of its absolute magnitude can be determined without knowing its distance nor the interstellar extinction.
In measuring star brightnesses, absolute magnitude, apparent magnitude, and distance are interrelated parameters—if two are known, the third can be determined. Since the Sun's luminosity is the standard, comparing these parameters with the Sun's apparent magnitude and distance is the easiest way to remember how to convert between them, although officially, zero point values are defined by the IAU.
The magnitude of a star, a unitless measure, is a logarithmic scale of observed visible brightness. The apparent magnitude is the observed visible brightness from Earth which depends on the distance of the object. The absolute magnitude is the apparent magnitude at a distance of , therefore the bolometric absolute magnitude is a logarithmic measure of the bolometric luminosity.
The difference in bolometric magnitude between two objects is related to their luminosity ratio according to:
where:
is the bolometric magnitude of the first object
is the bolometric magnitude of the second object.
is the first object's bolometric luminosity
is the second object's bolometric luminosity
The zero point of the absolute magnitude scale is actually defined as a fixed luminosity of . Therefore, the absolute magnitude can be calculated from a luminosity in watts:
where is the zero point luminosity
and the luminosity in watts can be calculated from an absolute magnitude (although absolute magnitudes are often not measured relative to an absolute flux):
| Physical sciences | Observational astronomy | null |
44883 | https://en.wikipedia.org/wiki/Welding | Welding | Welding is a fabrication process that joins materials, usually metals or thermoplastics, primarily by using high temperature to melt the parts together and allow them to cool, causing fusion. Common alternative methods include solvent welding (of thermoplastics) using chemicals to melt materials being bonded without heat, and solid-state welding processes which bond without melting, such as pressure, cold welding, and diffusion bonding.
Metal welding is distinct from lower temperature bonding techniques such as brazing and soldering, which do not melt the base metal (parent metal) and instead require flowing a filler metal to solidify their bonds.
In addition to melting the base metal in welding, a filler material is typically added to the joint to form a pool of molten material (the weld pool) that cools to form a joint that can be stronger than the base material. Welding also requires a form of shield to protect the filler metals or melted metals from being contaminated or oxidized.
Many different energy sources can be used for welding, including a gas flame (chemical), an electric arc (electrical), a laser, an electron beam, friction, and ultrasound. While often an industrial process, welding may be performed in many different environments, including in open air, under water, and in outer space. Welding is a hazardous undertaking and precautions are required to avoid burns, electric shock, vision damage, inhalation of poisonous gases and fumes, and exposure to intense ultraviolet radiation.
Until the end of the 19th century, the only welding process was forge welding, which blacksmiths had used for millennia to join iron and steel by heating and hammering. Arc welding and oxy-fuel welding were among the first processes to develop late in the century, and electric resistance welding followed soon after. Welding technology advanced quickly during the early 20th century, as world wars drove the demand for reliable and inexpensive joining methods. Following the wars, several modern welding techniques were developed, including manual methods like shielded metal arc welding, now one of the most popular welding methods, as well as semi-automatic and automatic processes such as gas metal arc welding, submerged arc welding, flux-cored arc welding and electroslag welding. Developments continued with the invention of laser beam welding, electron beam welding, magnetic pulse welding, and friction stir welding in the latter half of the century. Today, as the science continues to advance, robot welding is commonplace in industrial settings, and researchers continue to develop new welding methods and gain greater understanding of weld quality.
Etymology
The term weld is derived from the Middle English verb well (; plural/present tense: ) or welling (), meaning 'to heat' (to the maximum temperature possible); 'to bring to a boil'. The modern word was probably derived from the past-tense participle welled (), with the addition of d for this purpose being common in the Germanic languages of the Angles and Saxons. It was first recorded in English in 1590. A fourteenth century translation of the Christian Bible into English by John Wycliffe translates Isaiah 2:4 as "" (they shall beat together their swords into plowshares). In the 1590 version this was changed to "" (they shall weld together their swords into plowshares), suggesting this particular use of the word probably became popular in English sometime between these periods.
The Old English word for welding iron was ('to bring together') or ('to bring together hot').
The word is related to the Old Swedish word , meaning 'to boil', which could refer to joining metals, as in (literally 'to boil iron'). Sweden was a large exporter of iron during the Middle Ages, so the word may have entered English from the Swedish iron trade, or may have been imported with the thousands of Viking settlements that arrived in England before and during the Viking Age, as more than half of the most common English words in everyday use are Scandinavian in origin.
History
The history of joining metals goes back several millennia. The earliest examples of this come from the Bronze and Iron Ages in Europe and the Middle East. The ancient Greek historian Herodotus states in The Histories of the 5th century BC that Glaucus of Chios "was the man who single-handedly invented iron welding". Forge welding was used in the construction of the Iron pillar of Delhi, erected in Delhi, India about 310 AD and weighing 5.4 metric tons.
The Middle Ages brought advances in forge welding, in which blacksmiths pounded heated metal repeatedly until bonding occurred. In 1540, Vannoccio Biringuccio published De la pirotechnia, which includes descriptions of the forging operation. Renaissance craftsmen were skilled in the process, and the industry continued to grow during the following centuries.
In 1800, Sir Humphry Davy discovered the short-pulse electrical arc and presented his results in 1801. In 1802, Russian scientist Vasily Petrov created the continuous electric arc, and subsequently published "News of Galvanic-Voltaic Experiments" in 1803, in which he described experiments carried out in 1802. Of great importance in this work was the description of a stable arc discharge and the indication of its possible use for many applications, one being melting metals. In 1808, Davy, who was unaware of Petrov's work, rediscovered the continuous electric arc. In 1881–82 inventors Nikolai Benardos (Russian) and Stanisław Olszewski (Polish) created the first electric arc welding method known as carbon arc welding using carbon electrodes. The advances in arc welding continued with the invention of metal electrodes in the late 1800s by a Russian, Nikolai Slavyanov (1888), and an American, C. L. Coffin (1890). Around 1900, A. P. Strohmenger released a coated metal electrode in Britain, which gave a more stable arc. In 1905, Russian scientist Vladimir Mitkevich proposed using a three-phase electric arc for welding. Alternating current welding was invented by C. J. Holslag in 1919, but did not become popular for another decade.
Resistance welding was also developed during the final decades of the 19th century, with the first patents going to Elihu Thomson in 1885, who produced further advances over the next 15 years. Thermite welding was invented in 1893, and around that time another process, oxyfuel welding, became well established. Acetylene was discovered in 1836 by Edmund Davy, but its use was not practical in welding until about 1900, when a suitable torch was developed. At first, oxyfuel welding was one of the more popular welding methods due to its portability and relatively low cost. As the 20th century progressed, however, it fell out of favor for industrial applications. It was largely replaced with arc welding, as advances in metal coverings (known as flux) were made. Flux covering the electrode primarily shields the base material from impurities, but also stabilizes the arc and can add alloying components to the weld metal.
World War I caused a major surge in the use of welding, with the various military powers attempting to determine which of the several new welding processes would be best. The British primarily used arc welding, even constructing a ship, the "Fullagar" with an entirely welded hull. Arc welding was first applied to aircraft during the war as well, as some German airplane fuselages were constructed using the process. Also noteworthy is the first welded road bridge in the world, the Maurzyce Bridge in Poland (1928).
During the 1920s, significant advances were made in welding technology, including the introduction of automatic welding in 1920, in which electrode wire was fed continuously. Shielding gas became a subject receiving much attention, as scientists attempted to protect welds from the effects of oxygen and nitrogen in the atmosphere. Porosity and brittleness were the primary problems, and the solutions that developed included the use of hydrogen, argon, and helium as welding atmospheres. During the following decade, further advances allowed for the welding of reactive metals like aluminum and magnesium. This in conjunction with developments in automatic welding, alternating current, and fluxes fed a major expansion of arc welding during the 1930s and then during World War II. In 1930, the first all-welded merchant vessel, M/S Carolinian, was launched.
During the middle of the century, many new welding methods were invented. In 1930, Kyle Taylor was responsible for the release of stud welding, which soon became popular in shipbuilding and construction. Submerged arc welding was invented the same year and continues to be popular today. In 1932 a Russian, Konstantin Khrenov eventually implemented the first underwater electric arc welding. Gas tungsten arc welding, after decades of development, was finally perfected in 1941, and gas metal arc welding followed in 1948, allowing for fast welding of non-ferrous materials but requiring expensive shielding gases. Shielded metal arc welding was developed during the 1950s, using a flux-coated consumable electrode, and it quickly became the most popular metal arc welding process. In 1957, the flux-cored arc welding process debuted, in which the self-shielded wire electrode could be used with automatic equipment, resulting in greatly increased welding speeds, and that same year, plasma arc welding was invented by Robert Gage. Electroslag welding was introduced in 1958, and it was followed by its cousin, electrogas welding, in 1961. In 1953, the Soviet scientist N. F. Kazakov proposed the diffusion bonding method.
Other recent developments in welding include the 1958 breakthrough of electron beam welding, making deep and narrow welding possible through the concentrated heat source. Following the invention of the laser in 1960, laser beam welding debuted several decades later, and has proved to be especially useful in high-speed, automated welding. Magnetic pulse welding (MPW) has been industrially used since 1967. Friction stir welding was invented in 1991 by Wayne Thomas at The Welding Institute (TWI, UK) and found high-quality applications all over the world. All of these four new processes continue to be quite expensive due to the high cost of the necessary equipment, and this has limited their applications.
Processes
Welding joins two pieces of metal using heat, pressure, or both. The most common modern welding methods use heat sufficient to melt the base metals to be joined and the filler metal. This includes gas welding and all forms of arc welding. The area where the base and filler metals melt is called the weld pool or puddle. Most welding methods involve pushing the puddle along a joint to create a weld bead. Overlapping pieces of metal can be joined by forming the weld pool within a hole made in the topmost piece of base metal. This is called a plug weld. Overlapping base metals are commonly joined using electric resistance welding, a process that combines heat and pressure and does not require a filler metal. Solid-state welding processes join two pieces of metal using pressure.
Gas welding
The most common gas welding process is oxyfuel welding, also known as oxyacetylene welding. It is one of the oldest and most versatile welding processes, but in recent years it has become less popular in industrial applications. It is still widely used for welding pipes and tubes, as well as repair work.
The equipment is relatively inexpensive and simple, generally employing the combustion of acetylene in oxygen to produce a welding flame temperature of about 3100 °C (5600 °F). The flame, since it is less concentrated than an electric arc, causes slower weld cooling, which can lead to greater residual stresses and weld distortion, though it eases the welding of high alloy steels. A similar process, generally called oxyfuel cutting, is used to cut metals.
Arc welding
These processes use a welding power supply to create and maintain an electric arc between an electrode and the base material to melt metals at the welding point. They can use either direct current (DC) or alternating current (AC), and consumable or non-consumable electrodes. The welding region is sometimes protected by some type of inert or semi-inert gas, known as a shielding gas, and filler material is sometimes used as well.
Arc welding processes
One of the most common types of arc welding is shielded metal arc welding (SMAW); it is also known as manual metal arc welding (MMAW) or stick welding. Electric current is used to strike an arc between the base material and consumable electrode rod, which is made of filler material (typical steel) and is covered with a flux that protects the weld area from oxidation and contamination by producing carbon dioxide (CO2) gas during the welding process. The electrode core itself acts as filler material, making a separate filler unnecessary.
The process is versatile and can be performed with relatively inexpensive equipment, making it well suited to shop jobs and field work. An operator can become reasonably proficient with a modest amount of training and can achieve mastery with experience. Weld times are rather slow, since the consumable electrodes must be frequently replaced and because slag, the residue from the flux, must be chipped away after welding. Furthermore, the process is generally limited to welding ferrous materials, though special electrodes have made possible the welding of cast iron, stainless steel, aluminum, and other metals.
Gas metal arc welding (GMAW), also known as metal inert gas or MIG welding, is a semi-automatic or automatic process that uses a continuous wire feed as an electrode and an inert or semi-inert gas mixture to protect the weld from contamination. Since the electrode is continuous, welding speeds are greater for GMAW than for SMAW.
A related process, flux-cored arc welding (FCAW), uses similar equipment but uses wire consisting of a steel electrode surrounding a powder fill material. This cored wire is more expensive than the standard solid wire and can generate fumes and/or slag, but it permits even higher welding speed and greater metal penetration.
Gas tungsten arc welding (GTAW), or tungsten inert gas (TIG) welding, is a manual welding process that uses a non-consumable tungsten electrode, an inert or semi-inert gas mixture, and a separate filler material. Especially useful for welding thin materials, this method is characterized by a stable arc and high-quality welds, but it requires significant operator skill and can only be accomplished at relatively low speeds.
GTAW can be used on nearly all weldable metals, though it is most often applied to stainless steel and light metals. It is often used when quality welds are extremely important, such as in bicycle, aircraft and naval applications. A related process, plasma arc welding, also uses a tungsten electrode but uses plasma gas to make the arc. The arc is more concentrated than the GTAW arc, making transverse control more critical and thus generally restricting the technique to a mechanized process. Because of its stable current, the method can be used on a wider range of material thicknesses than can the GTAW process and it is much faster. It can be applied to all of the same materials as GTAW except magnesium, and automated welding of stainless steel is one important application of the process. A variation of the process is plasma cutting, an efficient steel cutting process.
Submerged arc welding (SAW) is a high-productivity welding method in which the arc is struck beneath a covering layer of flux. This increases arc quality since contaminants in the atmosphere are blocked by the flux. The slag that forms on the weld generally comes off by itself, and combined with the use of a continuous wire feed, the weld deposition rate is high. Working conditions are much improved over other arc welding processes, since the flux hides the arc and almost no smoke is produced. The process is commonly used in industry, especially for large products and in the manufacture of welded pressure vessels. Other arc welding processes include atomic hydrogen welding, electroslag welding (ESW), electrogas welding, and stud arc welding. ESW is a highly productive, single-pass welding process for thicker materials between 1 inch (25 mm) and 12 inches (300 mm) in a vertical or close to vertical position.
Arc welding power supplies
To supply the electrical power necessary for arc welding processes, a variety of different power supplies can be used. The most common welding power supplies are constant current power supplies and constant voltage power supplies. In arc welding, the length of the arc is directly related to the voltage, and the amount of heat input is related to the current. Constant current power supplies are most often used for manual welding processes such as gas tungsten arc welding and shielded metal arc welding, because they maintain a relatively constant current even as the voltage varies. This is important because in manual welding, it can be difficult to hold the electrode perfectly steady, and as a result, the arc length and thus voltage tend to fluctuate. Constant voltage power supplies hold the voltage constant and vary the current, and as a result, are most often used for automated welding processes such as gas metal arc welding, flux-cored arc welding, and submerged arc welding. In these processes, arc length is kept constant, since any fluctuation in the distance between the wire and the base material is quickly rectified by a large change in current. For example, if the wire and the base material get too close, the current will rapidly increase, which in turn causes the heat to increase and the tip of the wire to melt, returning it to its original separation distance.
The type of current used plays an important role in arc welding. Consumable electrode processes such as shielded metal arc welding and gas metal arc welding generally use direct current, but the electrode can be charged either positively or negatively. In welding, the positively charged anode will have a greater heat concentration, and as a result, changing the polarity of the electrode affects weld properties. If the electrode is positively charged, the base metal will be hotter, increasing weld penetration and welding speed. Alternatively, a negatively charged electrode results in more shallow welds. Non-consumable electrode processes, such as gas tungsten arc welding, can use either type of direct current, as well as alternating current. However, with direct current, because the electrode only creates the arc and does not provide filler material, a positively charged electrode causes shallow welds, while a negatively charged electrode makes deeper welds. Alternating current rapidly moves between these two, resulting in medium-penetration welds. One disadvantage of AC, the fact that the arc must be re-ignited after every zero crossings, has been addressed with the invention of special power units that produce a square wave pattern instead of the normal sine wave, making rapid zero crossings possible and minimizing the effects of the problem.
Resistance welding
Resistance welding involves the generation of heat by passing current through the resistance caused by the contact between two or more metal surfaces. Small pools of molten metal are formed at the weld area as high current (1,000–100,000 A) is passed through the metal. In general, resistance welding methods are efficient and cause little pollution, but their applications are somewhat limited and the equipment cost can be high.
Resistance spot welding is a popular method used to join overlapping metal sheets of up to 3 mm thick. Two electrodes are simultaneously used to clamp the metal sheets together and to pass current through the sheets. The advantages of the method include efficient energy use, limited workpiece deformation, high production rates, easy automation, and no required filler materials. Weld strength is significantly lower than with other welding methods, making the process suitable for only certain applications. It is used extensively in the automotive industry—ordinary cars can have several thousand spot welds made by industrial robots. A specialized process called shot welding, can be used to spot weld stainless steel.
Seam welding also relies on two electrodes to apply pressure and current to join metal sheets. However, instead of pointed electrodes, wheel-shaped electrodes roll along and often feed the workpiece, making it possible to make long continuous welds. In the past, this process was used in the manufacture of beverage cans, but now its uses are more limited. Other resistance welding methods include butt welding, flash welding, projection welding, and upset welding.
Energy beam welding
Energy beam welding methods, namely laser beam welding and electron beam welding, are relatively new processes that have become quite popular in high production applications. The two processes are quite similar, differing most notably in their source of power. Laser beam welding employs a highly focused laser beam, while electron beam welding is done in a vacuum and uses an electron beam. Both have a very high energy density, making deep weld penetration possible and minimizing the size of the weld area. Both processes are extremely fast, and are easily automated, making them highly productive. The primary disadvantages are their very high equipment costs (though these are decreasing) and a susceptibility to thermal cracking. Developments in this area include laser-hybrid welding, which uses principles from both laser beam welding and arc welding for even better weld properties, laser cladding, and x-ray welding.
Solid-state welding
Like forge welding (the earliest welding process discovered), some modern welding methods do not involve the melting of the materials being joined. One of the most popular, ultrasonic welding, is used to connect thin sheets or wires made of metal or thermoplastic by vibrating them at high frequency and under high pressure. The equipment and methods involved are similar to that of resistance welding, but instead of electric current, vibration provides energy input. When welding metals, the vibrations are introduced horizontally, and the materials are not melted; with plastics, which should have similar melting temperatures, vertically. Ultrasonic welding is commonly used for making electrical connections out of aluminum or copper, and it is also a very common polymer welding process.
Another common process, explosion welding, involves the joining of materials by pushing them together under extremely high pressure. The energy from the impact plasticizes the materials, forming a weld, even though only a limited amount of heat is generated. The process is commonly used for welding dissimilar materials, including bonding aluminum to carbon steel in ship hulls and stainless steel or titanium to carbon steel in petrochemical pressure vessels.
Other solid-state welding processes include friction welding (including friction stir welding and friction stir spot welding), magnetic pulse welding, co-extrusion welding, cold welding, diffusion bonding, exothermic welding, high frequency welding, hot pressure welding, induction welding, and roll bonding.
Geometry
Welds can be geometrically prepared in many different ways. The five basic types of weld joints are the butt joint, lap joint, corner joint, edge joint, and T-joint (a variant of this last is the cruciform joint). Other variations exist as well—for example, double-V preparation joints are characterized by the two pieces of material each tapering to a single center point at one-half their height. Single-U and double-U preparation joints are also fairly common—instead of having straight edges like the single-V and double-V preparation joints, they are curved, forming the shape of a U. Lap joints are also commonly more than two pieces thick—depending on the process used and the thickness of the material, many pieces can be welded together in a lap joint geometry.
Many welding processes require the use of a particular joint design; for example, resistance spot welding, laser beam welding, and electron beam welding are most frequently performed on lap joints. Other welding methods, like shielded metal arc welding, are extremely versatile and can weld virtually any type of joint. Some processes can also be used to make multipass welds, in which one weld is allowed to cool, and then another weld is performed on top of it. This allows for the welding of thick sections arranged in a single-V preparation joint, for example.
After welding, a number of distinct regions can be identified in the weld area. The weld itself is called the fusion zone—more specifically, it is where the filler metal was laid during the welding process. The properties of the fusion zone depend primarily on the filler metal used, and its compatibility with the base materials. It is surrounded by the heat-affected zone, the area that had its microstructure and properties altered by the weld. These properties depend on the base material's behavior when subjected to heat. The metal in this area is often weaker than both the base material and the fusion zone, and is also where residual stresses are found.
Quality
Many distinct factors influence the strength of welds and the material around them, including the welding method, the amount and concentration of energy input, the weldability of the base material, filler material, and flux material, the design of the joint, and the interactions between all these factors.
For example, the factor of welding position influences weld quality, that welding codes & specifications may require testing—both welding procedures and welders—using specified welding positions: 1G (flat), 2G (horizontal), 3G (vertical), 4G (overhead), 5G (horizontal fixed pipe), or 6G (inclined fixed pipe).
To test the quality of a weld, either destructive or nondestructive testing methods are commonly used to verify that welds are free of defects, have acceptable levels of residual stresses and distortion, and have acceptable heat-affected zone (HAZ) properties. Types of welding defects include cracks, distortion, gas inclusions (porosity), non-metallic inclusions, lack of fusion, incomplete penetration, lamellar tearing, and undercutting.
The metalworking industry has instituted codes and specifications to guide welders, weld inspectors, engineers, managers, and property owners in proper welding technique, design of welds, how to judge the quality of welding procedure specification, how to judge the skill of the person performing the weld, and how to ensure the quality of a welding job. Methods such as visual inspection, radiography, ultrasonic testing, phased-array ultrasonics, dye penetrant inspection, magnetic particle inspection, or industrial computed tomography can help with detection and analysis of certain defects.
Heat-affected zone
The heat-affected zone (HAZ) is a ring surrounding the weld in which the temperature of the welding process, combined with the stresses of uneven heating and cooling, alters the heat-treatment properties of the alloy. The effects of welding on the material surrounding the weld can be detrimental—depending on the materials used and the heat input of the welding process used, the HAZ can be of varying size and strength. The thermal diffusivity of the base material plays a large role—if the diffusivity is high, the material cooling rate is high and the HAZ is relatively small. Conversely, a low diffusivity leads to slower cooling and a larger HAZ. The amount of heat injected by the welding process plays an important role as well, as processes like oxyacetylene welding have an unconcentrated heat input and increase the size of the HAZ. Processes like laser beam welding give a highly concentrated, limited amount of heat, resulting in a small HAZ. Arc welding falls between these two extremes, with the individual processes varying somewhat in heat input. To calculate the heat input for arc welding procedures, the following formula can be used:
where Q = heat input (kJ/mm), V = voltage (V), I = current (A), and S = welding speed (mm/min). The efficiency is dependent on the welding process used, with shielded metal arc welding having a value of 0.75, gas metal arc welding and submerged arc welding, 0.9, and gas tungsten arc welding, 0.8. Methods of alleviating the stresses and brittleness created in the HAZ include stress relieving and tempering.
One major defect concerning the HAZ would be cracking at the toes , due to the rapid expansion (heating) and contraction (cooling) the material may not have the ability to withstand the stress and could cause cracking, one method the control these stress would be to control the heating and cooling rate, such as pre-heating and post- heating
Lifetime extension with after treatment methods
The durability and life of dynamically loaded, welded steel structures is determined in many cases by the welds, in particular the weld transitions. Through selective treatment of the transitions by grinding (abrasive cutting), shot peening, High-frequency impact treatment, Ultrasonic impact treatment, etc. the durability of many designs increases significantly.
Metallurgy
Most solids used are engineering materials consisting of crystalline solids in which the atoms or ions are arranged in a repetitive geometric pattern which is known as a lattice structure. The only exception is material that is made from glass which is a combination of a supercooled liquid and polymers which are aggregates of large organic molecules.
Crystalline solids cohesion is obtained by a metallic or chemical bond that is formed between the constituent atoms. Chemical bonds can be grouped into two types consisting of ionic and covalent. To form an ionic bond, either a valence or bonding electron separates from one atom and becomes attached to another atom to form oppositely charged ions. The bonding in the static position is when the ions occupy an equilibrium position where the resulting force between them is zero. When the ions are exerted in tension force, the inter-ionic spacing increases creating an electrostatic attractive force, while a repulsing force under compressive force between the atomic nuclei is dominant.
Covalent bonding takes place when one of the constituent atoms loses one or more electrons, with the other atom gaining the electrons, resulting in an electron cloud that is shared by the molecule as a whole. In both ionic and covalent bonding the location of the ions and electrons are constrained relative to each other, thereby resulting in the bond being characteristically brittle.
Metallic bonding can be classified as a type of covalent bonding for which the constituent atoms are of the same type and do not combine with one another to form a chemical bond. Atoms will lose an electron(s) forming an array of positive ions. These electrons are shared by the lattice which makes the electron cluster mobile, as the electrons are free to move as well as the ions. For this, it gives metals their relatively high thermal and electrical conductivity as well as being characteristically ductile.
Three of the most commonly used crystal lattice structures in metals are the body-centred cubic, face-centred cubic and close-packed hexagonal. Ferritic steel has a body-centred cubic structure and austenitic steel, non-ferrous metals like aluminium, copper and nickel have the face-centred cubic structure.
Ductility is an important factor in ensuring the integrity of structures by enabling them to sustain local stress concentrations without fracture. In addition, structures are required to be of an acceptable strength, which is related to a material's yield strength. In general, as the yield strength of a material increases, there is a corresponding reduction in fracture toughness.
A reduction in fracture toughness may also be attributed to the embrittlement effect of impurities, or for body-centred cubic metals, from a reduction in temperature. Metals and in particular steels have a transitional temperature range where above this range the metal has acceptable notch-ductility while below this range the material becomes brittle. Within the range, the materials behavior is unpredictable. The reduction in fracture toughness is accompanied by a change in the fracture appearance. When above the transition, the fracture is primarily due to micro-void coalescence, which results in the fracture appearing fibrous. When the temperatures falls the fracture will show signs of cleavage facets. These two appearances are visible by the naked eye. Brittle fracture in steel plates may appear as chevron markings under the microscope. These arrow-like ridges on the crack surface point towards the origin of the fracture.
Fracture toughness is measured using a notched and pre-cracked rectangular specimen, of which the dimensions are specified in standards, for example ASTM E23. There are other means of estimating or measuring fracture toughness by the following: The Charpy impact test per ASTM A370; The crack-tip opening displacement (CTOD) test per BS 7448–1; The J integral test per ASTM E1820; The Pellini drop-weight test per ASTM E208.
Unusual conditions
While many welding applications are done in controlled environments such as factories and repair shops, some welding processes are commonly used in a wide variety of conditions, such as open air, underwater, and vacuums (such as space). In open-air applications, such as construction and outdoors repair, shielded metal arc welding is the most common process. Processes that employ inert gases to protect the weld cannot be readily used in such situations, because unpredictable atmospheric movements can result in a faulty weld. Shielded metal arc welding is also often used in underwater welding in the construction and repair of ships, offshore platforms, and pipelines, but others, such as flux cored arc welding and gas tungsten arc welding, are also common. Welding in space is also possible—it was first attempted in 1969 by Russian cosmonauts during the Soyuz 6 mission, when they performed experiments to test shielded metal arc welding, plasma arc welding, and electron beam welding in a depressurized environment. Further testing of these methods was done in the following decades, and today researchers continue to develop methods for using other welding processes in space, such as laser beam welding, resistance welding, and friction welding. Advances in these areas may be useful for future endeavours similar to the construction of the International Space Station, which could rely on welding for joining in space the parts that were manufactured on Earth.
Safety issues
Welding can be dangerous and unhealthy if the proper precautions are not taken. However, using new technology and proper protection greatly reduces risks of injury and death associated with welding.
Since many common welding procedures involve an open electric arc or flame, the risk of burns and fire is significant; this is why it is classified as a hot work process. To prevent injury, welders wear personal protective equipment in the form of heavy leather gloves and protective long-sleeve jackets to avoid exposure to extreme heat and flames. Synthetic clothing such as polyester should not be worn since it may burn, causing injury. Additionally, the brightness of the weld area leads to a condition called arc eye or flash burns in which ultraviolet light causes inflammation of the cornea and can burn the retinas of the eyes. Goggles and welding helmets with dark UV-filtering face plates are worn to prevent this exposure. Since the 2000s, some helmets have included a face plate which instantly darkens upon exposure to the intense UV light. To protect bystanders, the welding area is often surrounded with translucent welding curtains. These curtains, made of a polyvinyl chloride plastic film, shield people outside the welding area from the UV light of the electric arc, but cannot replace the filter glass used in helmets. Depending on the type of material, welding varieties, and other factors, welding can produce over 100 dB(A) of noise. Long term or continuous exposure to higher decibels can lead to noise-induced hearing loss.
Welders are often exposed to dangerous gases and particulate matter. Processes like flux-cored arc welding and shielded metal arc welding produce smoke containing particles of various types of oxides. The size of the particles in question tends to influence the toxicity of the fumes, with smaller particles presenting a greater danger. This is because smaller particles have the ability to cross the blood–brain barrier. Fumes and gases, such as carbon dioxide, ozone, and fumes containing heavy metals, can be dangerous to welders lacking proper ventilation and training. Exposure to manganese welding fumes, for example, even at low levels (<0.2 mg/m3), may lead to neurological problems or to damage to the lungs, liver, kidneys, or central nervous system. Nano particles can become trapped in the alveolar macrophages of the lungs and induce pulmonary fibrosis. The use of compressed gases and flames in many, welding processes poses an explosion and fire risk. Some common precautions include limiting the amount of oxygen in the air, and keeping combustible materials away from the workplace.
Costs and trends
As an industrial process, the cost of welding plays a crucial role in manufacturing decisions. Many different variables affect the total cost, including equipment cost, labor cost, material cost, and energy cost. Depending on the process, equipment cost can vary, from inexpensive for methods like shielded metal arc welding and oxyfuel welding, to extremely expensive for methods like laser beam welding and electron beam welding. Because of their high cost, they are only used in high production operations. Similarly, because automation and robots increase equipment costs, they are only implemented when high production is necessary. Labor cost depends on the deposition rate (the rate of welding), the hourly wage, and the total operation time, including time spent fitting, welding, and handling the part. The cost of materials includes the cost of the base and filler material, and the cost of shielding gases. Finally, energy cost depends on arc time and welding power demand.
For manual welding methods, labor costs generally make up the vast majority of the total cost. As a result, many cost-saving measures are focused on minimizing operation time. To do this, welding procedures with high deposition rates can be selected, and weld parameters can be fine-tuned to increase welding speed. Mechanization and automation are often implemented to reduce labor costs, but this frequently increases the cost of equipment and creates additional setup time. Material costs tend to increase when special properties are necessary, and energy costs normally do not amount to more than several percent of the total welding cost.
In recent years, in order to minimize labor costs in high production manufacturing, industrial welding has become increasingly more automated, most notably with the use of robots in resistance spot welding (especially in the automotive industry) and in arc welding. In robot welding, mechanized devices both hold the material and perform the weld and at first, spot welding was its most common application, but robotic arc welding increases in popularity as technology advances. Other key areas of research and development include the welding of dissimilar materials (such as steel and aluminum, for example) and new welding processes, such as friction stir, magnetic pulse, conductive heat seam, and laser-hybrid welding. Furthermore, progress is desired in making more specialized methods like laser beam welding practical for more applications, such as in the aerospace and automotive industries. Researchers also hope to better understand the often unpredictable properties of welds, especially microstructure, residual stresses, and a weld's tendency to crack or deform.
The trend of accelerating the speed at which welds are performed in the steel erection industry comes at a risk to the integrity of the connection. Without proper fusion to the base materials provided by sufficient arc time on the weld, a project inspector cannot ensure the effective diameter of the puddle weld therefore he or she cannot guarantee the published load capacities unless they witness the actual installation. This method of puddle welding is common in the United States and Canada for attaching steel sheets to bar joist and structural steel members. Regional agencies are responsible for ensuring the proper installation of puddle welding on steel construction sites. Currently there is no standard or weld procedure which can ensure the published holding capacity of any unwitnessed connection, but this is under review by the American Welding Society.
Glass and plastic welding
Glasses and certain types of plastics are commonly welded materials. Unlike metals, which have a specific melting point, glasses and plastics have a melting range, called the glass transition. When heating the solid material past the glass-transition temperature (Tg) into this range, it will generally become softer and more pliable. When it crosses through the range, above the glass-melting temperature (Tm), it will become a very thick, sluggish, viscous liquid, slowly decreasing in viscosity as temperature increases. Typically, this viscous liquid will have very little surface tension compared to metals, becoming a sticky, taffy to honey-like consistency, so welding can usually take place by simply pressing two melted surfaces together. The two liquids will generally mix and join at first contact. Upon cooling through the glass transition, the welded piece will solidify as one solid piece of amorphous material.
Glass welding
Glass welding is a common practice during glassblowing. It is used very often in the construction of lighting, neon signs, flashtubes, scientific equipment, and the manufacture of dishes and other glassware. It is also used during glass casting for joining the halves of glass molds, making items such as bottles and jars. Welding glass is accomplished by heating the glass through the glass transition, turning it into a thick, formable, liquid mass. Heating is usually done with a gas or oxy-gas torch, or a furnace, because the temperatures for melting glass are often quite high. This temperature may vary, depending on the type of glass. For example, lead glass becomes a weldable liquid at around , and can be welded with a simple propane torch. On the other hand, quartz glass (fused silica) must be heated to over , but quickly loses its viscosity and formability if overheated, so an oxyhydrogen torch must be used. Sometimes a tube may be attached to the glass, allowing it to be blown into various shapes, such as bulbs, bottles, or tubes. When two pieces of liquid glass are pressed together, they will usually weld very readily. Welding a handle onto a pitcher can usually be done with relative ease. However, when welding a tube to another tube, a combination of blowing and suction, and pressing and pulling is used to ensure a good seal, to shape the glass, and to keep the surface tension from closing the tube in on itself. Sometimes a filler rod may be used, but usually not.
Because glass is very brittle in its solid state, it is often prone to cracking upon heating and cooling, especially if the heating and cooling are uneven. This is because the brittleness of glass does not allow for uneven thermal expansion. Glass that has been welded will usually need to be cooled very slowly and evenly through the glass transition, in a process called annealing, to relieve any internal stresses created by a temperature gradient.
There are many types of glass, and it is most common to weld using the same types. Different glasses often have different rates of thermal expansion, which can cause them to crack upon cooling when they contract differently. For instance, quartz has very low thermal expansion, while soda-lime glass has very high thermal expansion. When welding different glasses to each other, it is usually important to closely match their coefficients of thermal expansion, to ensure that cracking does not occur. Also, some glasses will simply not mix with others, so welding between certain types may not be possible.
Glass can also be welded to metals and ceramics, although with metals the process is usually more adhesion to the surface of the metal rather than a commingling of the two materials. However, certain glasses will typically bond only to certain metals. For example, lead glass bonds readily to copper or molybdenum, but not to aluminum. Tungsten electrodes are often used in lighting but will not bond to quartz glass, so the tungsten is often wetted with molten borosilicate glass, which bonds to both tungsten and quartz. However, care must be taken to ensure that all materials have similar coefficients of thermal expansion to prevent cracking both when the object cools and when it is heated again. Special alloys are often used for this purpose, ensuring that the coefficients of expansion match, and sometimes thin, metallic coatings may be applied to a metal to create a good bond with the glass.
Plastic welding
Plastics are generally divided into two categories, which are "thermosets" and "thermoplastics." A thermoset is a plastic in which a chemical reaction sets the molecular bonds after first forming the plastic, and then the bonds cannot be broken again without degrading the plastic. Thermosets cannot be melted, therefore, once a thermoset has set it is impossible to weld it. Examples of thermosets include epoxies, silicone, vulcanized rubber, polyester, and polyurethane.
Thermoplastics, by contrast, form long molecular chains, which are often coiled or intertwined, forming an amorphous structure without any long-range, crystalline order. Some thermoplastics may be fully amorphous, while others have a partially crystalline/partially amorphous structure. Both amorphous and semicrystalline thermoplastics have a glass transition, above which welding can occur, but semicrystallines also have a specific melting point which is above the glass transition. Above this melting point, the viscous liquid will become a free-flowing liquid (see rheological weldability for thermoplastics). Examples of thermoplastics include polyethylene, polypropylene, polystyrene, polyvinylchloride (PVC), and fluoroplastics like Teflon and Spectralon.
Welding thermoplastic with heat is very similar to welding glass. The plastic first must be cleaned and then heated through the glass transition, turning the weld-interface into a thick, viscous liquid. Two heated interfaces can then be pressed together, allowing the molecules to mix through intermolecular diffusion, joining them as one. Then the plastic is cooled through the glass transition, allowing the weld to solidify. A filler rod may often be used for certain types of joints. The main differences between welding glass and plastic are the types of heating methods, the much lower melting temperatures, and the fact that plastics will burn if overheated. Many different methods have been devised for heating plastic to a weldable temperature without burning it. Ovens or electric heating tools can be used to melt the plastic. Ultrasonic, laser, or friction heating are other methods. Resistive metals may be implanted in the plastic, which respond to induction heating. Some plastics will begin to burn at temperatures lower than their glass transition, so welding can be performed by blowing a heated, inert gas onto the plastic, melting it while, at the same time, shielding it from oxygen.
Solvent welding
Many thermoplastics can also be welded using chemical solvents. When placed in contact with the plastic, the solvent will begin to soften it, bringing the surface into a thick, liquid solution. When two melted surfaces are pressed together, the molecules in the solution mix, joining them as one. Because the solvent can permeate the plastic, the solvent evaporates out through the surface of the plastic, causing the weld to drop out of solution and solidify. A common use for solvent welding is for joining PVC (polyvinyl chloride) or ABS (acrylonitrile butadiene styrene) pipes during plumbing, or for welding styrene and polystyrene plastics in the construction of models. Solvent welding is especially effective on plastics like PVC which burn at or below their glass transition, but may be ineffective on plastics like Teflon or polyethylene that are resistant to chemical decomposition.
| Technology | Metallurgy | null |
44905 | https://en.wikipedia.org/wiki/Asthma | Asthma | Asthma is a common long-term inflammatory disease of the airways of the lungs. Asthma occurs when allergens, pollen, dust, or other particles, are inhaled into the lungs, causing the bronchioles to constrict and produce mucus, which then restricts oxygen flow to the alveoli. It is characterized by variable and recurring symptoms, reversible airflow obstruction, and easily triggered bronchospasms. Symptoms include episodes of wheezing, coughing, chest tightness, and shortness of breath. These may occur a few times a day or a few times per week. Depending on the person, asthma symptoms may become worse at night or with exercise.
Asthma is thought to be caused by a combination of genetic and environmental factors. Environmental factors include exposure to air pollution and allergens. Other potential triggers include medications such as aspirin and beta blockers. Diagnosis is usually based on the pattern of symptoms, response to therapy over time, and spirometry lung function testing. Asthma is classified according to the frequency of symptoms of forced expiratory volume in one second (FEV1), and peak expiratory flow rate. It may also be classified as atopic or non-atopic, where atopy refers to a predisposition toward developing a type 1 hypersensitivity reaction.
There is no known cure for asthma, but it can be controlled. Symptoms can be prevented by avoiding triggers, such as allergens and respiratory irritants, and suppressed with the use of inhaled corticosteroids. Long-acting beta agonists (LABA) or antileukotriene agents may be used in addition to inhaled corticosteroids if asthma symptoms remain uncontrolled. Treatment of rapidly worsening symptoms is usually with an inhaled short-acting beta2 agonist such as salbutamol and corticosteroids taken by mouth. In very severe cases, intravenous corticosteroids, magnesium sulfate, and hospitalization may be required.
In 2019 asthma affected approximately 262 million people and caused approximately 461,000 deaths. Most of the deaths occurred in the developing world. Asthma often begins in childhood, and the rates have increased significantly since the 1960s. Asthma was recognized as early as Ancient Egypt. The word asthma is from the Greek , , which means 'panting'.
Signs and symptoms
Asthma is characterized by recurrent episodes of wheezing, shortness of breath, chest tightness, and coughing. Sputum may be produced from the lung by coughing but is often hard to bring up. During recovery from an asthma attack (exacerbation), the sputum may appear pus-like due to high levels of white blood cells called eosinophils. Symptoms are usually worse at night and in the early morning or in response to exercise or cold air. Some people with asthma rarely experience symptoms, usually in response to triggers, whereas others may react frequently and readily and experience persistent symptoms.
Associated conditions
A number of other health conditions occur more frequently in people with asthma, including gastroesophageal reflux disease (GERD), rhinosinusitis, and obstructive sleep apnea. Psychological disorders are also more common, with anxiety disorders occurring in between 16 and 52% and mood disorders in 14–41%. It is not known whether asthma causes psychological problems or psychological problems lead to asthma. Current asthma, but not former asthma, is associated with increased all-cause mortality, heart disease mortality, and chronic lower respiratory tract disease mortality. Asthma, particularly severe asthma, is strongly associated with development of chronic obstructive pulmonary disease (COPD). Those with asthma, especially if it is poorly controlled, are at increased risk for radiocontrast reactions.
Cavities occur more often in people with asthma. This may be related to the effect of beta2-adrenergic agonists decreasing saliva. These medications may also increase the risk of dental erosions.
Causes
Asthma is caused by a combination of complex and incompletely understood environmental and genetic interactions. These influence both its severity and its responsiveness to treatment. It is believed that the recent increased rates of asthma are due to changing epigenetics (heritable factors other than those related to the DNA sequence) and a changing living environment. Asthma that starts before the age of 12 years old is more likely due to genetic influence, while onset after age 12 is more likely due to environmental influence.
Environmental
Many environmental factors have been associated with asthma's development and exacerbation, including allergens, air pollution, and other environmental chemicals. There are some substances that are known to cause asthma in exposed people and they are called asthmagens. Some common asthmagens include ammonia, latex, pesticides, solder and welding fumes, metal or wood dusts, spraying of isocyanate paint in vehicle repair, formaldehyde, glutaraldehyde, anhydrides, glues, dyes, metal working fluids, oil mists, moulds. Smoking during pregnancy and after delivery is associated with a greater risk of asthma-like symptoms. Low air quality from environmental factors such as traffic pollution or high ozone levels has been associated with both asthma development and increased asthma severity. Over half of cases in children in the United States occur in areas when air quality is below the EPA standards. Low air quality is more common in low-income and minority communities.
Exposure to indoor volatile organic compounds may be a trigger for asthma; formaldehyde exposure, for example, has a positive association. Phthalates in certain types of PVC are associated with asthma in both children and adults. While exposure to pesticides is linked to the development of asthma, a cause and effect relationship has yet to be established. A meta-analysis concluded gas stoves are a major risk factor for asthma, finding around one in eight cases in the U.S. could be attributed to these.
The majority of the evidence does not support a causal role between paracetamol (acetaminophen) or antibiotic use and asthma. A 2014 systematic review found that the association between paracetamol use and asthma disappeared when respiratory infections were taken into account. Maternal psychological stress during pregnancy is a risk factor for the child to develop asthma.
Asthma is associated with exposure to indoor allergens. Common indoor allergens include dust mites, cockroaches, animal dander (fragments of fur or feathers), and mould. Efforts to decrease dust mites have been found to be ineffective on symptoms in sensitized subjects. Weak evidence suggests that efforts to decrease mould by repairing buildings may help improve asthma symptoms in adults. Certain viral respiratory infections, such as respiratory syncytial virus and rhinovirus, may increase the risk of developing asthma when acquired as young children. Certain other infections, however, may decrease the risk.
Hygiene hypothesis
The hygiene hypothesis attempts to explain the increased rates of asthma worldwide as a direct and unintended result of reduced exposure, during childhood, to non-pathogenic bacteria and viruses. It has been proposed that the reduced exposure to bacteria and viruses is due, in part, to increased cleanliness and decreased family size in modern societies. Exposure to bacterial endotoxin in early childhood may prevent the development of asthma, but exposure at an older age may provoke bronchoconstriction. Evidence supporting the hygiene hypothesis includes lower rates of asthma on farms and in households with pets.
Use of antibiotics in early life has been linked to the development of asthma. Also, delivery via caesarean section is associated with an increased risk (estimated at 20–80%) of asthma – this increased risk is attributed to the lack of healthy bacterial colonization that the newborn would have acquired from passage through the birth canal. There is a link between asthma and the degree of affluence which may be related to the hygiene hypothesis as less affluent individuals often have more exposure to bacteria and viruses.
Genetic
Family history is a risk factor for asthma, with many different genes being implicated. If one identical twin is affected, the probability of the other having the disease is approximately 25%. By the end of 2005, 25 genes had been associated with asthma in six or more separate populations, including GSTM1, IL10, CTLA-4, SPINK5, LTC4S, IL4R and ADAM33, among others. Many of these genes are related to the immune system or modulating inflammation. Even among this list of genes supported by highly replicated studies, results have not been consistent among all populations tested. In 2006 over 100 genes were associated with asthma in one genetic association study alone; more continue to be found.
Some genetic variants may only cause asthma when they are combined with specific environmental exposures. An example is a specific single nucleotide polymorphism in the CD14 region and exposure to endotoxin (a bacterial product). Endotoxin exposure can come from several environmental sources including tobacco smoke, dogs, and farms. Risk for asthma, then, is determined by both a person's genetics and the level of endotoxin exposure.
Medical conditions
A triad of atopic eczema, allergic rhinitis and asthma is called atopy. The strongest risk factor for developing asthma is a history of atopic disease; with asthma occurring at a much greater rate in those who have either eczema or hay fever. Asthma has been associated with eosinophilic granulomatosis with polyangiitis (formerly known as Churg–Strauss syndrome), an autoimmune disease and vasculitis. Individuals with certain types of urticaria may also experience symptoms of asthma.
There is a correlation between obesity and the risk of asthma with both having increased in recent years. Several factors may be at play including decreased respiratory function due to a buildup of fat and the fact that adipose tissue leads to a pro-inflammatory state.
Beta blocker medications such as propranolol can trigger asthma in those who are susceptible. Cardioselective beta-blockers, however, appear safe in those with mild or moderate disease. Other medications that can cause problems in asthmatics are angiotensin-converting enzyme inhibitors, aspirin, and NSAIDs. Use of acid-suppressing medication (proton pump inhibitors and H2 blockers) during pregnancy is associated with an increased risk of asthma in the child.
Exacerbation
Some individuals will have stable asthma for weeks or months and then suddenly develop an episode of acute asthma. Different individuals react to various factors in different ways. Most individuals can develop severe exacerbation from a number of triggering agents.
Home factors that can lead to exacerbation of asthma include dust, animal dander (especially cat and dog hair), cockroach allergens and mold. Perfumes are a common cause of acute attacks in women and children. Both viral and bacterial infections of the upper respiratory tract can worsen the disease. Psychological stress may worsen symptoms – it is thought that stress alters the immune system and thus increases the airway inflammatory response to allergens and irritants.
Asthma exacerbations in school-aged children peak in autumn, shortly after children return to school. This might reflect a combination of factors, including poor treatment adherence, increased allergen and viral exposure, and altered immune tolerance. There is limited evidence to guide possible approaches to reducing autumn exacerbations, but while costly, seasonal omalizumab treatment from four to six weeks before school return may reduce autumn asthma exacerbations.
Pathophysiology
Asthma is the result of chronic inflammation of the conducting zone of the airways (most especially the bronchi and bronchioles), which subsequently results in increased contractability of the surrounding smooth muscles. This among other factors leads to bouts of narrowing of the airway and the classic symptoms of wheezing. The narrowing is typically reversible with or without treatment. Occasionally the airways themselves change. Typical changes in the airways include an increase in eosinophils and thickening of the lamina reticularis. Chronically the airways' smooth muscle may increase in size along with an increase in the numbers of mucous glands. Other cell types involved include T lymphocytes, macrophages, and neutrophils. There may also be involvement of other components of the immune system, including cytokines, chemokines, histamine, and leukotrienes among others.
Diagnosis
While asthma is a well-recognized condition, there is not one universal agreed-upon definition. It is defined by the Global Initiative for Asthma as "a chronic inflammatory disorder of the airways in which many cells and cellular elements play a role. The chronic inflammation is associated with airway hyper-responsiveness that leads to recurrent episodes of wheezing, breathlessness, chest tightness and coughing particularly at night or in the early morning. These episodes are usually associated with widespread but variable airflow obstruction within the lung that is often reversible either spontaneously or with treatment".
There is currently no precise test for the diagnosis, which is typically based on the pattern of symptoms and response to therapy over time. Asthma may be suspected if there is a history of recurrent wheezing, coughing or difficulty breathing and these symptoms occur or worsen due to exercise, viral infections, allergens or air pollution. Spirometry is then used to confirm the diagnosis. In children under the age of six the diagnosis is more difficult as they are too young for spirometry.
Spirometry
Spirometry is recommended to aid in diagnosis and management. It is the single best test for asthma. If the FEV1 measured by this technique improves more than 12% and increases by at least 200 millilitres following administration of a bronchodilator such as salbutamol, this is supportive of the diagnosis. It however may be normal in those with a history of mild asthma, not currently acting up. As caffeine is a bronchodilator in people with asthma, the use of caffeine before a lung function test may interfere with the results. Single-breath diffusing capacity can help differentiate asthma from COPD. It is reasonable to perform spirometry every one or two years to follow how well a person's asthma is controlled.
Others
The methacholine challenge involves the inhalation of increasing concentrations of a substance that causes airway narrowing in those predisposed. If negative it means that a person does not have asthma; if positive, however, it is not specific for the disease.
Other supportive evidence includes: a ≥20% difference in peak expiratory flow rate on at least three days in a week for at least two weeks, a ≥20% improvement of peak flow following treatment with either salbutamol, inhaled corticosteroids or prednisone, or a ≥20% decrease in peak flow following exposure to a trigger. Testing peak expiratory flow is more variable than spirometry, however, and thus not recommended for routine diagnosis. It may be useful for daily self-monitoring in those with moderate to severe disease and for checking the effectiveness of new medications. It may also be helpful in guiding treatment in those with acute exacerbations.
Classification
Asthma is clinically classified according to the frequency of symptoms, forced expiratory volume in one second (FEV1), and peak expiratory flow rate. Asthma may also be classified as atopic (extrinsic) or non-atopic (intrinsic), based on whether symptoms are precipitated by allergens (atopic) or not (non-atopic). While asthma is classified based on severity, at the moment there is no clear method for classifying different subgroups of asthma beyond this system. Finding ways to identify subgroups that respond well to different types of treatments is a current critical goal of asthma research. Recently, asthma has been classified based on whether it is associated with type 2 or non–type 2 inflammation. This approach to immunologic classification is driven by a developing understanding of the underlying immune processes and by the development of therapeutic approaches that target type 2 inflammation.
Although asthma is a chronic obstructive condition, it is not considered as a part of chronic obstructive pulmonary disease, as this term refers specifically to combinations of disease that are irreversible such as bronchiectasis and emphysema. Unlike these diseases, the airway obstruction in asthma is usually reversible; however, if left untreated, the chronic inflammation from asthma can lead the lungs to become irreversibly obstructed due to airway remodelling. In contrast to emphysema, asthma affects the bronchi, not the alveoli. The combination of asthma with a component of irreversible airway obstruction has been termed the asthma-chronic obstructive disease (COPD) overlap syndrome (ACOS). Compared to other people with "pure" asthma or COPD, people with ACOS exhibit increased morbidity, mortality and possibly more comorbidities.
Asthma exacerbation
An acute asthma exacerbation is commonly referred to as an asthma attack. The classic symptoms are shortness of breath, wheezing, and chest tightness. The wheezing is most often when breathing out. While these are the primary symptoms of asthma, some people present primarily with coughing, and in severe cases, air motion may be significantly impaired such that no wheezing is heard. In children, chest pain is often present.
Signs occurring during an asthma attack include the use of accessory muscles of respiration (sternocleidomastoid and scalene muscles of the neck), there may be a paradoxical pulse (a pulse that is weaker during inhalation and stronger during exhalation), and over-inflation of the chest. A blue colour of the skin and nails may occur from lack of oxygen.
In a mild exacerbation the peak expiratory flow rate (PEFR) is ≥200 L/min, or ≥50% of the predicted best. Moderate is defined as between 80 and 200 L/min, or 25% and 50% of the predicted best, while severe is defined as ≤ 80 L/min, or ≤25% of the predicted best.
Acute severe asthma, previously known as status asthmaticus, is an acute exacerbation of asthma that does not respond to standard treatments of bronchodilators and corticosteroids. Half of cases are due to infections with others caused by allergen, air pollution, or insufficient or inappropriate medication use.
Brittle asthma is a kind of asthma distinguishable by recurrent, severe attacks. Type 1 brittle asthma is a disease with wide peak flow variability, despite intense medication. Type 2 brittle asthma is background well-controlled asthma with sudden severe exacerbations.
Exercise-induced
Exercise can trigger bronchoconstriction both in people with or without asthma. It occurs in most people with asthma and up to 20% of people without asthma. Exercise-induced bronchoconstriction is common in professional athletes. The highest rates are among cyclists (up to 45%), swimmers, and cross-country skiers. While it may occur with any weather conditions, it is more common when it is dry and cold. Inhaled beta2 agonists do not appear to improve athletic performance among those without asthma; however, oral doses may improve endurance and strength.
Occupational
Asthma as a result of (or worsened by) workplace exposures is a commonly reported occupational disease. Many cases, however, are not reported or recognized as such. It is estimated that 5–25% of asthma cases in adults are work-related. A few hundred different agents have been implicated, with the most common being isocyanates, grain and wood dust, colophony, soldering flux, latex, animals, and aldehydes. The employment associated with the highest risk of problems include those who spray paint, bakers and those who process food, nurses, chemical workers, those who work with animals, welders, hairdressers and timber workers.
Aspirin-exacerbated respiratory disease
Aspirin-exacerbated respiratory disease (AERD), also known as aspirin-induced asthma, affects up to 9% of asthmatics. AERD consists of asthma, nasal polyps, sinus disease, and respiratory reactions to aspirin and other NSAID medications (such as ibuprofen and naproxen). People often also develop loss of smell and most experience respiratory reactions to alcohol.
Alcohol-induced asthma
Alcohol may worsen asthmatic symptoms in up to a third of people. This may be even more common in some ethnic groups such as the Japanese and those with aspirin-exacerbated respiratory disease. Other studies have found improvement in asthmatic symptoms from alcohol.
Non-atopic asthma
Non-atopic asthma, also known as intrinsic or non-allergic, makes up between 10 and 33% of cases. There is negative skin test to common inhalant allergens. Often it starts later in life, and women are more commonly affected than men. Usual treatments may not work as well. The concept that "non-atopic" is synonymous with "non-allergic" is called into question by epidemiological data that the prevalence of asthma is closely related to the serum IgE level standardized for age and sex (P<0.0001), indicating that asthma is almost always associated with some sort of IgE-related reaction and therefore has an allergic basis, although not all the allergic stimuli that cause asthma appear to have been included in the battery of aeroallergens studied (the "missing antigen(s)" hypothesis). For example, an updated systematic review and meta-analysis of population-attributable risk (PAR) of Chlamydia pneumoniae biomarkers in chronic asthma found that the PAR for C. pneumoniae-specific IgE was 47%.
Infectious asthma
Infectious asthma is an easily identified clinical presentation. When queried, asthma patients may report that their first asthma symptoms began after an acute lower respiratory tract illness. This type of history has been labelled the "infectious asthma" (IA) syndrome, or as "asthma associated with infection" (AAWI) to distinguish infection-associated asthma initiation from the well known association of respiratory infections with asthma exacerbations. Reported clinical prevalences of IA for adults range from around 40% in a primary care practice to 70% in a speciality practice treating mainly severe asthma patients. Additional information on the clinical prevalence of IA in adult-onset asthma is unavailable because clinicians are not trained to elicit this type of history routinely, and recollection in child-onset asthma is challenging. A population-based incident case-control study in a geographically defined area of Finland reported that 35.8% of new-onset asthma cases had experienced acute bronchitis or pneumonia in the year preceding asthma onset, representing a significantly higher risk compared to randomly selected controls (odds ratio 7.2, 95% confidence interval 5.2–10).
Phenotyping and endotyping
Asthma phenotyping and endotyping has emerged as a novel approach to asthma classification inspired by precision medicine which separates the clinical presentations of asthma, or asthma phenotypes, from their underlying causes, or asthma endotypes. The best-supported endotypic distinction is the type 2-high/type 2-low distinction. Classification based on type 2 inflammation is useful in predicting which patients will benefit from targeted biologic therapy.
Differential diagnosis
Many other conditions can cause symptoms similar to those of asthma. In children, symptoms may be due to other upper airway diseases such as allergic rhinitis and sinusitis, as well as other causes of airway obstruction including foreign body aspiration, tracheal stenosis, laryngotracheomalacia, vascular rings, enlarged lymph nodes or neck masses. Bronchiolitis and other viral infections may also produce wheezing. According to European Respiratory Society, it may not be suitable to label wheezing preschool children with the term asthma because there is lack of clinical data on inflammation in airways. In adults, COPD, congestive heart failure, airway masses, as well as drug-induced coughing due to ACE inhibitors may cause similar symptoms. In both populations vocal cord dysfunction may present similarly.
Chronic obstructive pulmonary disease can coexist with asthma and can occur as a complication of chronic asthma. After the age of 65, most people with obstructive airway disease will have asthma and COPD. In this setting, COPD can be differentiated by increased airway neutrophils, abnormally increased wall thickness, and increased smooth muscle in the bronchi. However, this level of investigation is not performed due to COPD and asthma sharing similar principles of management: corticosteroids, long-acting beta-agonists, and smoking cessation. It closely resembles asthma in symptoms, is correlated with more exposure to cigarette smoke, an older age, less symptom reversibility after bronchodilator administration, and decreased likelihood of family history of atopy.
Prevention
The evidence for the effectiveness of measures to prevent the development of asthma is weak. The World Health Organization recommends decreasing risk factors such as tobacco smoke, air pollution, chemical irritants including perfume, and the number of lower respiratory infections. Other efforts that show promise include: limiting smoke exposure in utero, breastfeeding, and increased exposure to daycare or large families, but none are well supported enough to be recommended for this indication.
Early pet exposure may be useful. Results from exposure to pets at other times are inconclusive and it is only recommended that pets be removed from the home if a person has allergic symptoms to said pet.
Dietary restrictions during pregnancy or when breastfeeding have not been found to be effective at preventing asthma in children and are not recommended. Omega-3 consumption, Mediterranean diet and antioxidants have been suggested by some studies to potentially help prevent crises but the evidence is still inconclusive.
Reducing or eliminating compounds known to sensitive people from the workplace may be effective. It is not clear if annual influenza vaccinations affect the risk of exacerbations. Immunization, however, is recommended by the World Health Organization. Smoking bans are effective in decreasing exacerbations of asthma.
Management
While there is no cure for asthma, symptoms can typically be improved. The most effective treatment for asthma is identifying triggers, such as cigarette smoke, pets or other allergens, and eliminating exposure to them. If trigger avoidance is insufficient, the use of medication is recommended. Pharmaceutical drugs are selected based on, among other things, the severity of illness and the frequency of symptoms. Specific medications for asthma are broadly classified into fast-acting and long-acting categories. The medications listed below have demonstrated efficacy in improving asthma symptoms; however, real world use-effectiveness is limited as around half of people with asthma worldwide remain sub-optimally controlled, even when treated. People with asthma may remain sub-optimally controlled either because optimum doses of asthma medications do not work (called "refractory" asthma) or because individuals are either unable (e.g. inability to afford treatment, poor inhaler technique) or unwilling (e.g., wish to avoid side effects of corticosteroids) to take optimum doses of prescribed asthma medications (called "difficult to treat" asthma). In practice, it is not possible to distinguish "refractory" from "difficult to treat" categories for patients who have never taken optimum doses of asthma medications. A related issue is that the asthma efficacy trials upon which the pharmacological treatment guidelines are based have systematically excluded the majority of people with asthma. For example, asthma efficacy treatment trials always exclude otherwise eligible people who smoke, and smoking diminishes the efficacy of inhaled corticosteroids, the mainstay of asthma control management.
Bronchodilators are recommended for short-term relief of symptoms. In those with occasional attacks, no other medication is needed. If mild persistent disease is present (more than two attacks a week), low-dose inhaled corticosteroids or alternatively, a leukotriene antagonist or a mast cell stabilizer by mouth is recommended. For those who have daily attacks, a higher dose of inhaled corticosteroids is used. In a moderate or severe exacerbation, corticosteroids by mouth are added to these treatments.
People with asthma have higher rates of anxiety, psychological stress, and depression. This is associated with poorer asthma control. Cognitive behavioural therapy may improve quality of life, asthma control, and anxiety levels in people with asthma.
Improving people's knowledge about asthma and using a written action plan has been identified as an important component of managing asthma. Providing educational sessions that include information specific to a person's culture is likely effective. More research is necessary to determine if increasing preparedness and knowledge of asthma among school staff and families using home-based and school interventions results in long term improvements in safety for children with asthma. School-based asthma self-management interventions, which attempt to improve knowledge of asthma, its triggers and the importance of regular practitioner review, may reduce hospital admissions and emergency department visits. These interventions may also reduce the number of days children experience asthma symptoms and may lead to small improvements in asthma-related quality of life. More research is necessary to determine if shared decision-making is helpful for managing adults with asthma or if a personalized asthma action plan is effective and necessary. Some people with asthma use pulse oximeters to monitor their own blood oxygen levels during an asthma attack. However, there is no evidence regarding the use in these instances.
Lifestyle modification
Avoidance of triggers is a key component of improving control and preventing attacks. The most common triggers include allergens, smoke (from tobacco or other sources), air pollution, nonselective beta-blockers, and sulfite-containing foods. Cigarette smoking and second-hand smoke (passive smoke) may reduce the effectiveness of medications such as corticosteroids. Laws that limit smoking decrease the number of people hospitalized for asthma. Dust mite control measures, including air filtration, chemicals to kill mites, vacuuming, mattress covers and other methods had no effect on asthma symptoms. There is insufficient evidence to suggest that dehumidifiers are helpful for controlling asthma.
Overall, exercise is beneficial in people with stable asthma. Yoga could provide small improvements in quality of life and symptoms in people with asthma. More research is necessary to determine how effective weight loss is in improving quality of life, the usage of health care services, and adverse effects for people of all ages with asthma.
Findings suggest that the Wim Hof Method may reduce inflammation in healthy and non-healthy participants as it increases epinephrine levels, causing an increase in interleukin-10 and a decrease in pro-inflammatory cytokines.
Medications
Medications used to treat asthma are divided into two general classes: quick-relief medications used to treat acute symptoms; and long-term control medications used to prevent further exacerbation. Antibiotics are generally not needed for sudden worsening of symptoms or for treating asthma at any time.
Medications for asthma exacerbations
Short-acting beta2-adrenoceptor agonists (SABAs), such as salbutamol (albuterol USAN) are the first-line treatment for asthma symptoms. They are recommended before exercise in those with exercise-induced symptoms.
Anticholinergic medications, such as ipratropium, provide additional benefit when used in combination with SABA in those with moderate or severe symptoms and may prevent hospitalizations. Anticholinergic bronchodilators can also be used if a person cannot tolerate a SABA. If a child requires admission to hospital additional ipratropium does not appear to help over a SABA. For children over 2 years old with acute asthma symptoms, inhaled anticholinergic medications taken alone is safe but is not as effective as inhaled SABA or SABA combined with inhaled anticholinergic medication. Adults who receive combined inhaled medications, which include short-acting anticholinergics and SABA, may be at risk for increased adverse effects such as experiencing a tremor, agitation, and heart beat palpitations compared to people who are treated with SABAs alone.
Older, less selective adrenergic agonists, such as inhaled epinephrine, have similar efficacy to SABAs. They are, however, not recommended due to concerns regarding excessive cardiac stimulation.
Corticosteroids can also help with the acute phase of an exacerbation because of their antiinflammatory properties. The benefit of systemic and oral corticosteroids is well established. Inhaled or nebulized corticosteroids can also be used. For adults and children who are in the hospital due to acute asthma, systemic (IV) corticosteroids improve symptoms. A short course of corticosteroids after an acute asthma exacerbation may help prevent relapses and reduce hospitalizations.
Other remedies, less established, are intravenous or nebulized magnesium sulfate and helium mixed with oxygen. Aminophylline could be used with caution as well.
Mechanical ventilation is the last resort in case of severe hypoxemia.
Intravenous administration of the drug aminophylline does not provide an improvement in bronchodilation when compared to standard inhaled beta2 agonist treatment. Aminophylline treatment is associated with more adverse effects compared to inhaled beta2 agonist treatment.
Long–term control
Corticosteroids are generally considered the most effective treatment available for long-term control. Inhaled forms are usually used except in the case of severe persistent disease, in which oral corticosteroids may be needed. Dosage depends on the severity of symptoms. High dosage and long-term use might lead to the appearance of common adverse effects which are growth delay, adrenal suppression, and osteoporosis. Continuous (daily) use of an inhaled corticosteroid, rather than its intermitted use, seems to provide better results in controlling asthma exacerbations. Commonly used corticosteroids are budesonide, fluticasone, mometasone and ciclesonide.
Long-acting beta-adrenoceptor agonists (LABA) such as salmeterol and formoterol can improve asthma control, at least in adults, when given in combination with inhaled corticosteroids. In children this benefit is uncertain. When used without steroids they increase the risk of severe side-effects, and with corticosteroids they may slightly increase the risk. Evidence suggests that for children who have persistent asthma, a treatment regime that includes LABA added to inhaled corticosteroids may improve lung function but does not reduce the amount of serious exacerbations. Children who require LABA as part of their asthma treatment may need to go to the hospital more frequently.
Leukotriene receptor antagonists (anti-leukotriene agents such as montelukast and zafirlukast) may be used in addition to inhaled corticosteroids, typically also in conjunction with a LABA. For adults or adolescents who have persistent asthma that is not controlled very well, the addition of anti-leukotriene agents along with daily inhaled corticosteriods improves lung function and reduces the risk of moderate and severe asthma exacerbations. Anti-leukotriene agents may be effective alone for adolescents and adults; however, there is no clear research suggesting which people with asthma would benefit from anti-leukotriene receptor alone. In those under five years of age, anti-leukotriene agents were the preferred add-on therapy after inhaled corticosteroids. A 2013 Cochrane systematic review concluded that anti-leukotriene agents appear to be of little benefit when added to inhaled steroids for treating children. A similar class of drugs, 5-LOX inhibitors, may be used as an alternative in the chronic treatment of mild to moderate asthma among older children and adults. there is one medication in this family known as zileuton.
Mast cell stabilizers (such as cromolyn sodium) are safe alternatives to corticosteroids but not preferred because they have to be administered frequently.
Oral theophyllines are sometimes used for controlling chronic asthma, but their used is minimized due to side effects.
Omalizumab, a monoclonal antibody against IgE, is a novel way to lessen exacerbations by decreasing the levels of circulating IgE that play a significant role at allergic asthma.
Anticholinergic medications such as ipratropium bromide have not been shown to be beneficial for treating chronic asthma in children over 2 years old, and are not suggested for routine treatment of chronic asthma in adults.
There is no strong evidence to recommend chloroquine medication as a replacement for taking corticosteroids by mouth (for those who are not able to tolerate inhaled steroids). Methotrexate is not suggested as a replacement for taking corticosteriods by mouth ("steroid-sparing") due to the adverse effects associated with taking methotrexate and the minimal relief provided for asthma symptoms.
Macrolide antibiotics, particularly the azalide macrolide azithromycin, are a recently added Global Initiative for Asthma (GINA)-recommended treatment option for both eosinophilic and non-eosinophilic severe, refractory asthma based on azithromycin's efficacy in reducing moderate and severe exacerbations combined. Azithromycin's mechanism of action is not established, and could involve pathogen- and/or host-directed anti-inflammatory activities. Limited clinical observations suggest that some patients with new-onset asthma and with "difficult-to-treat" asthma (including those with the asthma-COPD overlap syndrome – ACOS) may respond dramatically to azithromycin. However, these groups of asthma patients have not been studied in randomized treatment trials and patient selection needs to be carefully individualized.
A 2024 study indicates that commonly used diabetes medications may lower asthma attacks by up to 70%. The research examined metformin and GLP-1 drugs such as Ozempic (semaglutide), Mounjaro (tirzepatide), and Saxenda (liraglutide). Among nearly 13,000 participants with both diabetes and asthma, metformin reduced the risk of asthma attacks by 30%, with an additional 40% reduction when combined with a GLP-1 drug.
For children with asthma which is well-controlled on combination therapy of inhaled corticosteroids (ICS) and long-acting beta2-agonists (LABA), the benefits and harms of stopping LABA and stepping down to ICS-only therapy are uncertain. In adults who have stable asthma while they are taking a combination of LABA and inhaled corticosteroids (ICS), stopping LABA may increase the risk of asthma exacerbations that require treatment with corticosteroids by mouth. Stopping LABA probably makes little or no important difference to asthma control or asthma-related quality of life. Whether or not stopping LABA increases the risk of serious adverse events or exacerbations requiring an emergency department visit or hospitalization is uncertain.
Delivery methods
Medications are typically provided as metered-dose inhalers (MDIs) in combination with an inhaler spacer or as a dry powder inhaler. The spacer is a plastic cylinder that mixes the medication with air, making it easier to receive a full dose of the drug. A nebulizer may also be used. Nebulizers and spacers are equally effective in those with mild to moderate symptoms. However, insufficient evidence is available to determine whether a difference exists in those with severe disease. For delivering short-acting beta-agonists in acute asthma in children, spacers may have advantages compared to nebulisers, but children with life-threatening asthma have not been studied. There is no strong evidence for the use of intravenous LABA for adults or children who have acute asthma. There is insufficient evidence to directly compare the effectiveness of a metered-dose inhaler attached to a homemade spacer compared to commercially available spacer for treating children with asthma.
Adverse effects
Long-term use of inhaled corticosteroids at conventional doses carries a minor risk of adverse effects. Risks include thrush, the development of cataracts, and a slightly slowed rate of growth. Rinsing the mouth after the use of inhaled steroids can decrease the risk of thrush. Higher doses of inhaled steroids may result in lower bone mineral density.
Others
Inflammation in the lungs can be estimated by the level of exhaled nitric oxide. The use of exhaled nitric oxide levels (FeNO) to guide asthma medication dosing may have small benefits for preventing asthma attacks but the potential benefits are not strong enough for this approach to be universally recommended as a method to guide asthma therapy in adults or children.
When asthma is unresponsive to usual medications, other options are available for both emergency management and prevention of flareups. Additional options include:
Humidified oxygen to alleviate hypoxia if saturations fall below 92%.
Corticosteroids by mouth, with five days of prednisone being the same two days of dexamethasone. One review recommended a seven-day course of steroids.
Magnesium sulfate intravenous treatment increases bronchodilation when used in addition to other treatment in moderate severe acute asthma attacks. In adults intravenous treatment results in a reduction of hospital admissions. Low levels of evidence suggest that inhaled (nebulized) magnesium sulfate may have a small benefit for treating acute asthma in adults. Overall, high-quality evidence do not indicate a large benefit for combining magnesium sulfate with standard inhaled treatments for adults with asthma.
Heliox, a mixture of helium and oxygen, may also be considered in severe unresponsive cases.
Intravenous salbutamol is not supported by available evidence and is thus used only in extreme cases.
Methylxanthines (such as theophylline) were once widely used, but do not add significantly to the effects of inhaled beta-agonists. Their use in acute exacerbations is controversial.
The dissociative anaesthetic ketamine is theoretically useful if intubation and mechanical ventilation is needed in people who are approaching respiratory arrest; however, there is no evidence from clinical trials to support this. A 2012 Cochrane review found no significant benefit from the use of ketamine in severe acute asthma in children.
For those with severe persistent asthma not controlled by inhaled corticosteroids and LABAs, bronchial thermoplasty may be an option. It involves the delivery of controlled thermal energy to the airway wall during a series of bronchoscopies. While it may increase exacerbation frequency in the first few months it appears to decrease the subsequent rate. Effects beyond one year are unknown.
Monoclonal antibody injections such as mepolizumab, dupilumab, or omalizumab may be useful in those with poorly controlled atopic asthma. However, these medications are expensive and their use is therefore reserved for those with severe symptoms to achieve cost-effectiveness. Monoclonal antibodies targeting interleukin-5 (IL-5) or its receptor (IL-5R), including mepolizumab, reslizumab or benralizumab, in addition to standard care in severe asthma is effective in reducing the rate of asthma exacerbations. There is limited evidence for improved health-related quality of life and lung function.
Evidence suggests that sublingual immunotherapy in those with both allergic rhinitis and asthma improve outcomes.
It is unclear if non-invasive positive pressure ventilation in children is of use as it has not been sufficiently studied.
Adherence to asthma treatments
Staying with a treatment approach for preventing asthma exacerbations can be challenging, especially if the person is required to take medicine or treatments daily. Reasons for low adherence range from a conscious decision to not follow the suggested medical treatment regime for various reasons including avoiding potential side effects, misinformation, or other beliefs about the medication. Problems accessing the treatment and problems administering the treatment effectively can also result in lower adherence. Various approaches have been undertaken to try and improve adherence to treatments to help people prevent serious asthma exacerbations including digital interventions.
Alternative medicine
Many people with asthma, like those with other chronic disorders, use alternative treatments; surveys show that roughly 50% use some form of unconventional therapy. There is little data to support the effectiveness of most of these therapies.
Evidence is insufficient to support the usage of vitamin C or vitamin E for controlling asthma. There is tentative support for use of vitamin C in exercise induced bronchospasm. Fish oil dietary supplements (marine n-3 fatty acids) and reducing dietary sodium do not appear to help improve asthma control. In people with mild to moderate asthma, treatment with vitamin D supplementation or its hydroxylated metabolites does not reduce acute exacerbations or improve control. There is no strong evidence to suggest that vitamin D supplements improve day-to-day asthma symptoms or a person's lung function. There is no strong evidence to suggest that adults with asthma should avoid foods that contain monosodium glutamate (MSG). There have not been enough high-quality studies performed to determine if children with asthma should avoid eating food that contains MSG.
Acupuncture is not recommended for the treatment as there is insufficient evidence to support its use. Air ionizers show no evidence that they improve asthma symptoms or benefit lung function; this applied equally to positive and negative ion generators. Manual therapies, including osteopathic, chiropractic, physiotherapeutic and respiratory therapeutic manoeuvres, have insufficient evidence to support their use in treating asthma. Pulmonary rehabilitation, however, may improve quality of life and functional exercise capacity when compared to usual care for adults with asthma. The Buteyko breathing technique for controlling hyperventilation may result in a reduction in medication use; however, the technique does not have any effect on lung function. Thus an expert panel felt that evidence was insufficient to support its use. There is no clear evidence that breathing exercises are effective for treating children with asthma.
Prognosis
The prognosis for asthma is generally good, especially for children with mild disease. Mortality has decreased over the last few decades due to better recognition and improvement in care. In 2010 the death rate was 170 per million for males and 90 per million for females. Rates vary between countries by 100-fold.
Globally it causes moderate or severe disability in 19.4 million people (16 million of which are in low and middle income countries). Of asthma diagnosed during childhood, half of cases will no longer carry the diagnosis after a decade. Airway remodelling is observed, but it is unknown whether these represent harmful or beneficial changes. More recent data find that severe asthma can result in airway remodelling and the "asthma with chronic obstructive pulmonary disease syndrome (ACOS)" that has a poor prognosis. Early treatment with corticosteroids seems to prevent or ameliorates a decline in lung function. Asthma in children also has negative effects on quality of life of their parents.
Epidemiology
In 2019, approximately 262 million people worldwide were affected by asthma and approximately 461,000 people died from the disease. Rates vary between countries with prevalences between 1 and 18%. It is more common in developed than developing countries. One thus sees lower rates in Asia, Eastern Europe and Africa. Within developed countries it is more common in those who are economically disadvantaged while in contrast in developing countries it is more common in the affluent. The reason for these differences is not well known. Low- and middle-income countries make up more than 80% of the mortality.
While asthma is twice as common in boys as girls, severe asthma occurs at equal rates. In contrast adult women have a higher rate of asthma than men and it is more common in the young than the old. In 2010, children with asthma experienced over 900,000 emergency department visits, making it the most common reason for admission to the hospital following an emergency department visit in the US in 2011.
Global rates of asthma have increased significantly between the 1960s and 2008 with it being recognized as a major public health problem since the 1970s. Rates of asthma have plateaued in the developed world since the mid-1990s with recent increases primarily in the developing world. Asthma affects approximately 7% of the population of the United States and 5% of people in the United Kingdom. Canada, Australia and New Zealand have rates of about 14–15%.
The average death rate from 2011 to 2015 from asthma in the UK was about 50% higher than the average for the European Union and had increased by about 5% in that time. Children are more likely see a physician due to asthma symptoms after school starts in September.
Population-based epidemiological studies describe temporal associations between acute respiratory illnesses, asthma, and development of severe asthma with irreversible airflow limitation (known as the asthma-chronic obstructive pulmonary disease "overlap" syndrome, or ACOS). Additional prospective population-based data indicate that ACOS seems to represent a form of severe asthma, characterized by more frequent hospitalizations, and to be the result of early-onset asthma that has progressed to fixed airflow obstruction.
Economics
From 2000 to 2010, the average cost per asthma-related hospital stay in the United States for children remained relatively stable at about $3,600, whereas the average cost per asthma-related hospital stay for adults increased from $5,200 to $6,600. In 2010, Medicaid was the most frequent primary payer among children and adults aged 18–44 years in the United States; private insurance was the second most frequent payer. Among both children and adults in the lowest income communities in the United States there is a higher rate of hospital stays for asthma in 2010 than those in the highest income communities.
History
Asthma was recognized in ancient Egypt and was treated by drinking an incense mixture known as kyphi. It was officially named as a specific respiratory problem by Hippocrates circa 450 BC, with the Greek word for "panting" forming the basis of our modern name. In 200 BC it was believed to be at least partly related to the emotions. In the 12th century the Jewish physician-philosopher Maimonides wrote a treatise on asthma in Arabic, based partly on Arabic sources, in which he discussed the symptoms, proposed various dietary and other means of treatment, and emphasized the importance of climate and clean air. Traditional Chinese medicine also offered medication for asthma, as indicated by a surviving 14th-century manuscript curated by the Wellcome Foundation.
In 1873, one of the first papers in modern medicine on the subject tried to explain the pathophysiology of the disease while one in 1872, concluded that asthma can be cured by rubbing the chest with chloroform liniment. Medical treatment in 1880 included the use of intravenous doses of a drug called pilocarpine.
In 1886, F. H. Bosworth theorized a connection between asthma and hay fever.
At the beginning of the 20th century, the focus was the avoidance of allergens as well as selective beta-2 adrenoceptor agonists were used as treatment strategies.
Epinephrine was first referred to in the treatment of asthma in 1905. Oral corticosteroids began to be used for the condition in 1950. The use of a pressurized metered-dose inhaler was developed in the mid-1950s for the administration of adrenaline and isoproterenol and was later used as a beta2-adrenergic agonist.
Inhaled corticosteroids and selective short-acting beta agonists came into wide use in the 1960s.
A well-documented case in the 19th century was that of young Theodore Roosevelt (1858–1919). At that time there was no effective treatment. Roosevelt's youth was in large part shaped by his poor health, partly related to his asthma. He experienced recurring nighttime asthma attacks that felt as if he was being smothered to death, terrifying the boy and his parents.
During the 1930s to 1950s, asthma was known as one of the "holy seven" psychosomatic illnesses. Its cause was considered to be psychological, with treatment often based on psychoanalysis and other talking cures. As these psychoanalysts interpreted the asthmatic wheeze as the suppressed cry of the child for its mother, they considered the treatment of depression to be especially important for individuals with asthma.
In January 2021, an appeal court in France overturned a deportation order against a 40-year-old Bangladeshi man, who was a patient of asthma. His lawyers had argued that the dangerous levels of pollution in Bangladesh could possibly lead to worsening of his health condition, or even premature death.
| Biology and health sciences | Illness and injury | null |
44928 | https://en.wikipedia.org/wiki/Mahogany | Mahogany | Mahogany is a straight-grained, reddish-brown timber of three tropical hardwood species of the genus Swietenia, indigenous to the Americas and part of the pantropical chinaberry family, Meliaceae. Mahogany is used commercially for a wide variety of goods, due to its coloring and durable nature. It is naturally found within the Americas, but has also been imported to plantations across Asia and Oceania. The mahogany trade may have begun as early as the 16th century and flourished in the 17th and 18th centuries. In certain countries, mahogany is considered an invasive species.
Mahogany is wood from any of three tree species: Honduran or big-leaf mahogany (Swietenia macrophylla), West Indian or Cuban mahogany (Swietenia mahagoni), and Swietenia humilis. Honduran mahogany is the most widespread and the only genuine mahogany species commercially grown today. Mahogany is a valuable lumber used for paneling, furniture, boats, musical instruments, and other items. The United States is the leading importer of mahogany, while Peru is the largest exporter. Mahogany is the national tree of the Dominican Republic and Belize.
Swietenia species have been introduced in various countries outside the Americas since the 1800s, with many plantings becoming naturalized forests. All species of Swietenia are now listed by CITES and protected due to concerns over illegal logging and mismanagement. Mahogany species can crossbreed when they grow in proximity, and the hybrid between S. mahagoni and S. macrophylla is widely planted for timber.
The history of the American mahogany trade dates back to the 17th century when the wood was first noticed by Europeans during the Spanish colonization of the Americas. Mahogany became more popular in the 18th century when the British Parliament removed import duties on timber from British possessions, leading to increased exports to Europe and North America. Throughout the 18th and 19th centuries, mahogany from various regions was imported into Europe and North America, with Britain being the largest consumer.
By the late 19th century, African mahogany began to dominate the market, and by the early 20th century, the supply of American mahogany became scarcer. In response to concerns about the sustainability of mahogany, several species have been placed on CITES Appendices to regulate the trade.
Mahogany is known for its straight, fine grain and durability, making it a popular choice for fine furniture, boat construction, and musical instruments. However, the over-harvesting of mahogany and environmental concerns have led to a decrease in its use.
Etymology
The etymology of mahogany is uncertain and a subject of debate. The term first appeared in John Ogilby's "America" (1671), referring to a "curious and rich wood" from Jamaica. Initial mentions of the mahogany tree (as opposed to wood) date to 1731, with its first detailed description in 1743, attributed to Swietenia mahagoni by Kemp Malone in 1940. Malone suggested that mahogany originated as a generic term for 'wood' in a native Bahamian language. F. Bruce Lamb disagreed, pointing out that the Arawak language's word for wood is caoba. Lamb identified a West African origin for the word in the Yoruba oganwo, collectively m'oganwo (meaning one which is the tallest or most high) used for the Khaya genus of trees, whose timber is today called African mahogany. Lamb proposes that Yoruba and Igbo people brought to Jamaica as slaves identified the local trees of the Swietenia genus as m'oganwo, which developed into the Portuguese term mogano, which first appeared in print as the name of a river in 1661, before finally developing into the English mahogany in Jamaica between 1655 and 1670.
Malone criticized this etymology, arguing that the proposed metamorphosis from the Yoruba m'oganwo to the Portuguese mogano to the English mahogany was a logical and linguistic stretch relying on the conversion of the singular oganwo to the collective m'oganwo, which Malone finds unlikely considering the tree's generally solitary nature. He also argues that Lamb's earliest identified use of the Portuguese mog(a)no, which is for a river that Lamb asserts must have been so named for the mahogany oganwo trees on its banks, could just as well have been named for any tall tree, since oganwo only means tall. Lamb, in turn, criticized Malone's methodology and perceived bias, and maintained that there is no evidence for mahogany as a generic word.
Description
Mahogany is a commercially important lumber prized for its beauty, durability, and color, and used for paneling and to make furniture, boats, musical instruments and other items. The leading importer of mahogany is the United States, followed by Britain; while the largest exporter today is Peru, which surpassed Brazil after that country banned mahogany exports in 2001. It is estimated that some 80 or 90 percent of Peruvian mahogany exported to the United States is illegally harvested, with the economic cost of illegal logging in Peru placed conservatively at $40–70 million USD annually. It was estimated that in 2000, some 57,000 mahogany trees were harvested to supply the U.S. furniture trade alone.
Mahogany is the national tree of the Dominican Republic and Belize. A mahogany tree with two woodcutters bearing an axe and a paddle also appears on the Belizean national coat of arms, under the national motto, , Latin for "under the shade I flourish."
The specific density of mahogany is 0.55. Mahogany, African: (500–850 kg/m3); Mahogany, Cuban: 660 kg/m3; Mahogany, Honduras: 650 kg/m3; Mahogany, Spanish: 850 kg/m3.
Species
The three species are:
Honduran or big-leaf mahogany (Swietenia macrophylla), with a range from Mexico to southern Amazonia in Brazil, the most widespread species of mahogany and the only genuine mahogany species commercially grown today. Illegal logging of S. macrophylla, and its highly destructive environmental effects, led to the species' placement in 2003 on Appendix II of Convention on International Trade in Endangered Species (CITES), the first time that a high-volume, high-value tree was listed on Appendix II.
West Indian or Cuban mahogany (Swietenia mahagoni), native to southern Florida and the Caribbean, formerly dominant in the mahogany trade, but not in widespread commercial use since World War II.
Swietenia humilis, a small and often twisted mahogany tree limited to seasonally dry forests in Pacific Central America that is of limited commercial utility. Some botanists believe that S. humilis is a mere variant of S. macrophylla.
Other species
While only the three Swietenia species are classified officially as "genuine mahogany", the Federal Trade Commission allows certain species of trees other than Swietenia to be sold as "mahoganies" in the U.S. timber trade. This is due to the long-standing usage of the terms. But it must be prefixed with another descriptor, and they are not allowed to be sold under the name "mahogany" alone.
Two names are allowed. The first is "African mahogany" for the five species of the genus Khaya (which also belong to the mahogany family), namely: K. anthotheca, K. grandifoliola, K. ivorensis, K. madagascariensis, and K. senegalensis. All of them are native to native to Africa and Madagascar. The second is the name "Philippine mahogany" for seven species (all native to the Philippines) in the genus Shorea and Parashorea (which are unrelated dipterocarps, more commonly known as "lauan" or "meranti"), namely:S. polysperma, S. negrosensis, S. contorta, S. ovata, S. almon, S. palosapis, and P. malaanonan. The timber from both "African mahoganies" and "Philippine mahoganies" as defined by the FTC, are very close in terms of appearance and properties to true mahoganies. No other species are allowed to be sold in the United States under the name "mahogany", aside from the three Swietenia species and the aforementioned exceptions.
Within the mahogany family, other closely-related members of other genera which also resemble mahoganies in terms of appearance and properties are also sometimes known as "mahoganies", though they can not be sold as such in the US timber trade. This includes some members of the genus Toona, namely: "Philippine mahogany" (Toona calantas, different from the above usage); "Indian mahogany" (Toona ciliata); "Chinese mahogany" (Toona sinensis); and Indonesian mahogany (Toona sureni);. However members of this genus are more usually known as "toons" or "red cedars." They have similar properties to true mahoganies but differ in appearance. Other species in the same family sometimes known as "mahoganies" include "Indian mahogany" (Chukrasia velutina, different from T. ciliata); "sipo mahogany" (Entandrophragma utile); "sapele mahogany" (Entandrophragma cylindricum); "royal mahogany" (Carapa guianensis); "white mahogany" (Turraeanthus africanus); "New Zealand mahogany" (Dysoxylum spectabile); "pink mahogany" (Guarea spp.); and "demerara mahogany" (Carapa guianensis).
Multiple other unrelated species are also known as "mahogany". These include the aforementioned Shorea species which does actually come close to true mahogany in terms of appearance and properties. But it also includes other species which do not resemble true mahogany at all and have very different wood properties, like the "Santos mahogany" (Myroxylon balsamum), "mountain mahogany" (Cercocarpus spp.), and "swamp mahogany" (Eucalyptus robusta).
Distribution
The natural distribution of these species within the Americas is geographically distinct. S. mahagoni grows on the West Indian islands as far north as the Bahamas, the Florida Keys and parts of Florida; S. humilis grows in the dry regions of the Pacific coast of Central America from south-western Mexico to Costa Rica; S. macrophylla grows in Central America from Yucatan southwards and into South America, extending as far as Peru, Bolivia and extreme western Brazil. In the 20th century various botanists attempted to further define S. macrophylla in South America as a new species, such as S. candollei Pittier and S. tessmannii Harms., but many authorities consider these spurious. According to Record and Hess, all of the mahogany of continental North and South America can be considered as one botanical species, Swietenia macrophylla King.
Both major species of Swietenia were introduced in several countries outside of the Americas during the 1800s and early 1900s using seeds from South America and the Caribbean. Many of these plantings became naturalized forests over time.
India had both S. macrophylla and S. mahagoni introduced in 1865 using seeds from West Indies. Both eventually became naturalized forests. Bangladesh had Honduran S. macrophylla introduced in 1872 and as with India it became naturalized in some areas. S. mahagoni and S. macrophylla were introduced in Indonesia in 1870 using seeds from India. S. macrophylla was included in plantation forests planted in Indonesia from the 1920s to the 1940s. Philippines had S. macrophylla introduced in 1907 and in 1913 as well as S. mahagoni in 1911, 1913, 1914, 1920 and 1922. Planting resumed in the late 1980s. It was planted with many other exotic tree species for the purpose of reforestation. S. macrophylla was planted in Sri Lanka in 1897 but it was left unmanaged until the 1950s when reforestation efforts initiated by the Sri Lankan government led to plantations being consciously developed. In the early 1900s S. mahagoni was planted on the islands of O'ahu and Maui in Hawaii but was neglected and became naturalized forests. Additionally, S. macrophylla was planted in 1922 on O'ahu and is now naturalized. Fiji had S. macrophylla introduced originally in 1911 as an ornamental species using seeds from Honduras and Belize. Fiji has become a major producer of mahogany in the 21st Century due to a robust plantation program spanning over 50 years. Harvesting began in 2003.
History
The name mahogany was initially associated only with those islands in the West Indies under British control (French colonists used the term acajou, while in the Spanish territories it was called caoba). The origin of the name is uncertain, but it could be a corruption of 'm'oganwo', the name used by the Yoruba and Ibo people of West Africa to describe trees of the genus Khaya, which is closely related to Swietenia. When transported to Jamaica as slaves, they gave the same name to the similar trees they saw there. Though this interpretation has been disputed, no one has suggested a more plausible origin. The indigenous Arawak name for the tree is not known. In 1671 the word mahogany appeared in print for the first time, in John Ogilby's America. Among botanists and naturalists, however, the tree was considered a type of cedar, and in 1759 was classified by Carl Linnaeus (1707–1778) as Cedrela mahagoni. The following year it was assigned to a new genus by Nicholas Joseph Jacquin (1727–1817), and named Swietenia mahagoni.
Until the 19th century all of the mahogany was regarded as one species, although varying in quality and character according to soil and climate. In 1836 the German botanist Joseph Gerhard Zuccarini (1797–1848) identified a second species while working on specimens collected on the Pacific coast of Mexico, and named it Swietenia humilis. In 1886 a third species, Swietenia macrophylla, was named by Sir George King (1840–1909) after studying specimens of Honduras mahogany planted in the Botanic Gardens in Calcutta, India.
Today, all species of Swietenia grown in their native locations are listed by CITES, and are therefore protected. After S. mahogani and S. macrophylla were added to CITES appendixes in 1992 and 1995 respectively international conservation programs began in earnest aided by a 1993 World Bank report entitled "Tropical Hardwood Marketing Strategies for Southeast Asia". Efforts to repopulate mahogany largely failed in its native locations due to attacks from the shoot borer Hypsipyla grandella and similarly failed in Africa due to the attacks by the equivalent Hypsipyla robusta. After so many years of mismanagement and illegal logging, Swietenia also suffered from genetic loss thus mutating and weakening the seeds. Additionally erosion in its native locations meant seeds could no longer even be planted. However, both species grew well in Asia and Asia Pacific due to the absence of these shoot borers and absence of other limitations. Plantation management progressed throughout the 1990s and 2000s in Asia and the South Pacific. Global supply of genuine mahogany has been increasing from these plantations, notably Fiji, and Philippines. For Swietenia macrophylla, the trees in these plantations are still relatively young compared to the trees being harvested from old growth forests in South America. Thus, the illegal trade of bigleaf mahogany continues apace.
History of American mahogany trade
In the 17th century, the buccaneer Alexandre Exquemelin recorded the use of mahogany or Caoba (Cedrela being the Spanish name) on Hispaniola for making canoes: "The Indians make these canoes without the use of any iron instruments, by only burning the trees at the bottom near the root, and afterwards governing the fire with such industry that nothing is burnt more than what they would have..."
The wood first came to the notice of Europeans with the beginning of Spanish colonisation in the Americas. A cross in the Cathedral at Santo Domingo, bearing the date 1514, is said to be mahogany, and Philip II of Spain apparently used the wood for the interior joinery of the palace El Escorial, begun in 1584. However, caoba, as the Taino Natives called the wood, was principally reserved for shipbuilding, and it was declared a royal monopoly at Havana in 1622. Hence very little of the mahogany growing in Spanish controlled territory found its way to Europe.
After the French established a colony in Saint Domingue (now Haiti), some mahogany from that island probably found its way to France, where joiners in the port cities of Saint-Malo, Nantes, La Rochelle and Bordeaux used the wood to a limited extent from about 1700. On the English-controlled islands, especially Jamaica and the Bahamas, mahogany was abundant but not exported in any quantity before 1700.
18th century
While the trade in mahogany from the Spanish and French territories in America remained moribund for most of the 18th century, this was not true for those islands under British control. The British Parliament passed an act of Parliament, the Naval Stores Act 1721 (8 Geo. 1. c. 12), which removed all import duties from timber imported into Britain from British possessions in the Americas. This immediately stimulated the trade in West Indian timbers including, most importantly, mahogany. Importations of mahogany into England (and excluding those to Scotland, which were recorded separately) reached 525 tons per annum by 1740, 3,688 tons by 1750, and more than 30,000 tons in 1788, the peak year of the 18th century trade.
At the same time, the Naval Stores Act 1721 had the effect of substantially increasing exports of mahogany from the West Indies to the British colonies in North America. Although initially regarded as a joinery wood, mahogany rapidly became the timber of choice for makers of high quality furniture in both the British Isles and the 13 colonies of North America.
Until the 1760s over 90 per cent of the mahogany imported into Britain came from Jamaica. Some of this was re-exported to continental Europe, but most was used by British furniture makers. Quantities of Jamaican mahogany also went to the North American colonies, but most of the wood used in American furniture came from the Bahamas. This was sometimes called Providence wood, after the main port of the islands, but more often madera or madeira, which was the West Indian name for mahogany.
In addition to Jamaica and Bahamas, all the British controlled islands exported some mahogany at various times, but the quantities were not large. The most significant third source was Black River and adjacent areas on the Mosquito Coast (now Republic of Honduras), from where quantities of mahogany were shipped from the 1740s onwards. This mahogany was known as 'Rattan mahogany', after the island of Ruatan, which was the main offshore entrepot for the British settlers in the area.
At the end of the Seven Years' War (1756–63), the mahogany trade began to change significantly. During the occupation of Havana by British forces between August 1762 and July 1763, quantities of Cuban or Havanna mahogany were sent to Britain, and after the city was restored to Spain in 1763, Cuba continued to export small quantities, mostly to ports on the north coast of Jamaica, from where it went to Britain. However, this mahogany was regarded as inferior to the Jamaican variety, and the trade remained sporadic until the 19th century.
Another variety new to the market was Hispaniola mahogany, also called 'Spanish' and 'St Domingo' mahogany. This was the result of the Free Port Act 1766 (6 Geo. 3. c. 49), which opened Kingston and other designated Jamaican ports to foreign vessels for the first time. The object was primarily to encourage importations of cotton from French plantations in Saint Domingue, but quantities of high quality mahogany were also shipped. These were then forwarded to Britain, where they entered the market in the late 1760s.
In terms of quantity, the most significant new addition to the mahogany trade was Honduras mahogany, also called 'baywood', after the Bay of Honduras. British settlers had been active in southern Yucatan since the beginning of the 18th century, despite the opposition of the Spanish, who claimed sovereignty over all of Central America.
Their main occupation was cutting logwood, a dyewood in high demand in Europe. The center of their activity and the primary point of export was Belize. Under Article XVII of the Treaty of Paris (1763), British cutters were for the first time given the right to cut logwood in Yucatan unmolested, within agreed limits. Such was the enthusiasm of the cutters that within a few years the European market was glutted, and the price of logwood collapsed.
However, the price of mahogany was still high after the war, and so the cutters turned to cutting mahogany. The first Honduras mahogany arrived in Kingston, Jamaica, in November 1763, and the first shipments arrived in Britain the following year.
By the 1790s most of the viable stocks of mahogany in Jamaica had been cut, and the market was divided between two principal sources or types of mahogany. Honduras mahogany was relatively cheap, plentiful, but rarely of the best quality. Hispaniola (also called Spanish or Santo Domingo) mahogany was the wood of choice for high quality work.
Data are lacking, but it is likely that the newly independent United States now received a good proportion of its mahogany from Cuba. In the last quarter of the 18th century France began to use mahogany more widely; they had ample supplies of high quality wood from Saint Domingue. The rest of Europe, where the wood was increasingly fashionable, obtained most of their wood from Britain.
Recent history
The French Revolution of 1789 and the wars that followed radically changed the mahogany trade, primarily due to the progressive collapse of the French and Spanish colonial empires, which allowed British traders into areas previously closed to them. Saint Domingue became the independent republic of Haiti, and from 1808, Spanish controlled Santo Domingo and Cuba were both open to British vessels for the first time.
From the 1820s mahogany from all these areas was imported into Europe and North America, with the majority of them going to Britain. In Central America British loggers moved northwest towards Mexico and south into Guatemala. Other areas of Central America as far south as Panama also began to be exploited.
The most important new development was the beginning of large scale logging in Mexico from the 1860s. Most mahogany was cut in the province of Tabasco and exported from a number of ports on the Gulf of Campeche, from Vera Cruz eastwards to Campeche and Sisal. By the end of the 19th century there was scarcely any part of Central America within reach of the coast untouched by logging, and activity also extended into Colombia, Venezuela, Peru and Brazil.
Trade in American mahogany probably reached a peak in the last quarter of the 19th century. Figures are not available for all countries, but Britain alone imported more than 80,000 tons in 1875. This figure was not matched again. From the 1880s, African mahogany (Khaya spp.), a related genus, began to be exported in increasing quantities from West Africa, and by the early 20th century it dominated the market.
In 1907 the total of mahogany from all sources imported into Europe was 159,830 tons, of which 121,743 tons were from West Africa. By this time mahogany from Cuba, Haiti and other West Indian sources had become increasingly difficult to obtain in commercial sizes, and by the late 20th century Central American and even South American mahogany was heading in a similar direction. In 1975 S. humilis was placed on CITES Appendix II (a list of species that would be in danger of extinction without strict regulation) followed by S. mahagoni in 1992. The most abundant species, S. macrophylla, was placed on Appendix III in 1995 and moved to Appendix II in 2003.
Uses
Mahogany has a straight, fine, and even grain, and is relatively free of voids and pockets. Its reddish-brown color darkens over time, and displays a reddish sheen when polished. It has excellent workability, and is very durable. Historically, the tree's girth allowed for wide boards from traditional mahogany species. These properties make it a favorable wood for crafting cabinets and furniture.
Much of the first-quality furniture made in the American colonies from the mid 18th century was made of mahogany, when the wood first became available to American craftsmen. Mahogany is still widely used for fine furniture; however, the rarity of Cuban mahogany, the over-harvesting of Honduras and Brazilian mahogany, and the protests by indigenous peoples and environmental organizations from the 1980s into the 2000s, have diminished their use. Recent mahogany production from Mexico and Fiji has a lighter color and density than South American production from the early 20th century.
Mahogany also resists wood rot, making it attractive in boat construction and outdoor decking. It is a tonewood, often used for musical instruments, particularly the backs, sides and necks of acoustic guitars, electric guitar bodies, and drum shells because of its ability to produce a very deep, warm tone compared to other commonly used woods, such as maple, alder, ash (Fraxinus) or spruce. Guitars featuring mahogany in their construction include many acoustic guitars from Martin, Taylor, and Gibson, and Gibson electric guitars such as the Les Paul and SG. In the 1930s Gibson used the wood to make banjo necks as well.
Mahogany as an invasive species
In the Philippines, environmentalists are calling for an end to the planting of mahogany because of its negative impact on the environment and wildlife, including possibly causing soil acidification and no net benefit to wildlife.
| Biology and health sciences | Sapindales | null |
317350 | https://en.wikipedia.org/wiki/Common%20quail | Common quail | The common quail (Coturnix coturnix), or European quail, is a small ground-nesting game bird in the pheasant family Phasianidae. It is mainly migratory, breeding in the western Palearctic and wintering in Africa and southern India.
With its characteristic call of three repeated chirps (repeated three times in quick succession), this species of quail is more often heard than seen. It is widespread in Europe and North Africa, and is categorised by the IUCN as "least concern". It should not be confused with the Japanese quail (Coturnix japonica), native to Asia, which, although visually similar, has a call that is very distinct from that of the common quail. Like the Japanese quail, common quails are sometimes kept as poultry.
Taxonomy
The common quail was formally described by the Swedish naturalist Carl Linnaeus in 1758 in the tenth edition of his Systema Naturae under the binomial name Tetrao coturnix. The specific epithet coturnix is the Latin word for the common quail. This species is now placed in the genus Coturnix that was introduced in 1764 by the French naturalist François Alexandre Pierre de Garsault. The common quail was formerly considered to be conspecific with the Japanese quail (Coturnix japonica). The ranges of the two species meet in Mongolia and near Lake Baikal without apparent interbreeding and in captivity the offspring of crosses show reduced fertility. The Japanese quail is therefore now treated as a separate species.
Five subspecies are recognised:
C. c. coturnix (Linnaeus, 1758) – breeding in Europe and northwest Africa to Mongolia and north India, wintering in Africa and central, south India
C. c. conturbans Hartert, 1917 – Azores
C. c. inopinata Hartert, 1917 – Cape Verde Islands
C. c. africana Temminck & Schlegel, 1848 – sub-Saharan Africa and the three islands
C. c. erlangeri Zedlitz, 1912 – east and northeast Africa
Description
The common quail is a small compact gallinaceous bird in length with a wingspan of . The weight is . It is greatest before migration at the end of the breeding season. The female is generally slightly heavier than the male. It is streaked brown with a white eyestripe, and, in the male, a white chin. As befits its migratory nature, it has long wings, unlike the typically short-winged gamebirds. According to Online Etymology Dictionary, "small migratory game bird of the Old World, late 14c. (early 14c. as a surname, Quayle), from Old French quaille (Modern French caille), perhaps via Medieval Latin quaccula (source also of Provençal calha, Italian quaglia, Portuguese calha, Old Spanish coalla), or directly from a Germanic source (compare Dutch kwakkel, Old High German quahtala, German Wachtel, Old English wihtel), imitative of the bird's cry. Or the English word might have come up indigenously from Proto-Germanic."
Distribution and habitat
This is a terrestrial species, feeding on seeds and insects on the ground. It is notoriously difficult to see, keeping hidden in crops, and reluctant to fly, preferring to creep away instead. Even when flushed, it keeps low and soon drops back into cover. Often the only indication of its presence is the distinctive "wet-my-lips" repetitive song of the male. The call is uttered mostly in the mornings, evenings and sometimes at night. It is a strongly migratory bird, unlike most game birds.
The common quail has been introduced onto the island of Mauritius on several occasions but has failed to establish itself and is now probably extinct.
Behaviour and ecology
Breeding
Males generally arrive in the breeding area before the females. In northern Europe laying begins from the middle of May, and with repeat laying can continue to the end of August. The female forms a shallow scrape in the ground in diameter which is sparsely lined with vegetation. The eggs are laid at 24-hour intervals to form a clutch of between 8 and 13 eggs. These have an off-white to creamy yellow background with dark brown spots or blotches. Their average dimensions are with a weight of . The eggs are incubated by the female alone beginning after all the eggs are laid. The eggs hatch synchronously after 17–20 days. The young are precocial and shortly after hatching leave the nest and can feed themselves. They are cared for by the female who broods them while they are small. The young fledge when around 19 days of age but stay in the family group for 30–50 days. They generally first breed when one year old and only have a single brood.
Relationship to humans
The common quail is heavily hunted as game on passage through the Mediterranean area. Very large numbers are caught in nets along the Mediterranean coast of Egypt. It is estimated that in 2012, during the autumn migration, 3.4 million birds were caught in northern Sinai and perhaps as many as 12.9 million in the whole of Egypt.
This species over recent years has seen an increase in its propagation in the United States and Europe. However, most of this increase is with hobbyists. It is declining in parts of its range such as Ireland.
In 1537, Queen Jane Seymour, third wife of Henry VIII, then pregnant with the future King Edward VI, developed an insatiable craving for quail, and courtiers and diplomats abroad were ordered to find sufficient supplies for the Queen.
Poisoning
If common quails have eaten certain plants, although which plants is still in debate, the meat from quail can be poisonous, with one in four who consume poisonous flesh becoming ill with coturnism, which is characterized by muscle soreness, and which may lead to kidney failure.
In culture
In the Bible, the Book of Numbers chapter 11 describes a story of a huge mass of quails that were blown by a wind and were taken as meat by the Israelites in the wilderness.
Gallery
| Biology and health sciences | Galliformes | null |
317380 | https://en.wikipedia.org/wiki/Menthol | Menthol | Menthol is an organic compound, specifically a monoterpenoid, that occurs naturally in the oils of several plants in the mint family, such as corn mint and peppermint. It is a white or clear waxy crystalline substance that is solid at room temperature and melts slightly above. The main form of menthol occurring in nature is (−)-menthol, which is assigned the (1R,2S,5R) configuration.
For many people, menthol produces a cooling sensation when inhaled, eaten, or applied to the skin, and mint plants have been used for centuries for topical pain relief and as a food flavoring. Menthol has local anesthetic and counterirritant qualities, and it is widely used to relieve minor throat irritation. Menthol also acts as a weak κ-opioid receptor agonist.
Structure
Natural menthol exists as one pure stereoisomer, nearly always the (1R,2S,5R) form (bottom left corner of the diagram below). The eight possible stereoisomers are:
In the natural compound, the isopropyl group is in the trans orientation to both the methyl and hydroxyl groups. Thus, it can be drawn in any of the ways shown:
The (+)- and (−)-enantiomers of menthol are the most stable among these based on their cyclohexane conformations. With the ring itself in a chair conformation, all three bulky groups can orient in equatorial positions.
The two crystal forms for racemic menthol have melting points of 28 °C and 38 °C. Pure (−)-menthol has four crystal forms, of which the most stable is the α form, the familiar broad needles.
Biological properties
Menthol's ability to chemically trigger the cold-sensitive TRPM8 receptors in the skin is responsible for the well-known cooling sensation it provokes when inhaled, eaten, or applied to the skin. In this sense, it is similar to capsaicin, the chemical responsible for the spiciness of hot chilis (which stimulates heat sensors, also without causing an actual change in temperature).
Menthol's analgesic properties are mediated through a selective activation of κ-opioid receptors. Menthol blocks calcium channels and voltage-sensitive sodium channels, reducing neural activity that may stimulate muscles.
Some studies show that menthol acts as a GABAA receptor positive allosteric modulator and increases GABAergic transmission in PAG neurons. Menthol has anesthetic properties similar to, though less potent than, propofol because it interacts with the same sites on the GABAA receptor. Menthol may also enhance the activity of glycine receptors and negatively modulate 5-HT3 receptors and nAChRs.
Menthol is widely used in dental care as a topical antibacterial agent, effective against several types of streptococci and lactobacilli. Menthol also lowers blood pressure and antagonizes vasoconstriction through TRPM8 activation.
Occurrence
Mentha arvensis (wild mint) is the primary species of mint used to make natural menthol crystals and natural menthol flakes. This species is primarily grown in the Uttar Pradesh region in India.
Menthol occurs naturally in peppermint oil (along with a little menthone, the ester menthyl acetate and other compounds), obtained from Mentha × piperita (peppermint). Japanese menthol also contains a small percentage of the 1-epimer neomenthol.
Biosynthesis
The biosynthesis of menthol has been investigated in Mentha × piperita and the enzymes involved in have been identified and characterized. It begins with the synthesis of the terpene limonene, followed by hydroxylation, and then several reduction and isomerization steps.
More specifically, the biosynthesis of (−)-menthol takes place in the secretory gland cells of the peppermint plant. The steps of the biosynthetic pathway are as follows:
Geranyl diphosphate synthase (GPPS) first catalyzes the reaction of IPP and DMAPP into geranyl diphosphate.
(−)-limonene synthase (LS) catalyzes the cyclization of geranyl diphosphate to (−)-limonene.
(−)-Limonene-3-hydroxylase (L3OH), using O2 and then nicotinamide adenine dinucleotide phosphate (NADPH) catalyzes the allylic hydroxylation of (−)-limonene at the 3 position to (−)-trans-isopiperitenol.
(−)-trans-Isopiperitenol dehydrogenase (iPD) further oxidizes the hydroxyl group on the 3 position using NAD+ to make (−)-isopiperitenone.
(−)-Isopiperitenone reductase (iPR) then reduces the double bond between carbons 1 and 2 using NADPH to form (+)-cis-isopulegone.
(+)-cis-Isopulegone isomerase (iPI) then isomerizes the remaining double bond to form (+)-pulegone.
(+)-Pulegone reductase (PR) reduces this double bond using NADPH to form (−)-menthone.
(−)-Menthone reductase (MR) then reduces the carbonyl group using NADPH to form (−)-menthol.
Production
Natural menthol is obtained by freezing peppermint oil. The resultant crystals of menthol are then separated by filtration.
Total world production of menthol in 1998 was 12,000 tonnes of which 2,500 tonnes was synthetic. In 2005, the annual production of synthetic menthol was almost double. Prices are in the $10–20/kg range with peaks in the $40/kg region but have reached as high as $100/kg. In 1985, it was estimated that China produced most of the world's supply of natural menthol, although it appears that India has pushed China into second place.
Menthol is manufactured as a single enantiomer (94% e.e.) on the scale of 3,000 tonnes per year by Takasago International Corporation. The process involves an asymmetric synthesis developed by a team led by Ryōji Noyori, who won the 2001 Nobel Prize for Chemistry in recognition of his work on this process:
The process begins by forming an allylic amine from myrcene, which undergoes asymmetric isomerisation in the presence of a BINAP rhodium complex to give (after hydrolysis) enantiomerically pure R-citronellal. This is cyclised by a carbonyl-ene-reaction initiated by zinc bromide to , which is then hydrogenated to give pure (1R,2S,5R)-menthol.
Another commercial process is the Haarmann–Reimer process (after the company Haarmann & Reimer, now part of Symrise) This process starts from m-cresol which is alkylated with propene to thymol. This compound is hydrogenated in the next step. Racemic menthol is isolated by fractional distillation. The enantiomers are separated by chiral resolution in reaction with methyl benzoate, selective crystallisation followed by hydrolysis.
Racemic menthol can also be formed by hydrogenation of thymol, menthone, or pulegone. In both cases with further processing (crystallizative entrainment resolution of the menthyl benzoate conglomerate) it is possible to concentrate the L-enantiomer, however this tends to be less efficient, although the higher processing costs may be offset by lower raw material costs. A further advantage of this process is that D-menthol becomes inexpensively available for use as a chiral auxiliary, along with the more usual L-antipode.
Applications
Menthol is included in many products, and for a variety of reasons.
Cosmetic
In some beauty products such as hair conditioners, based on natural ingredients (e.g., St. Ives).
Medical
As an antipruritic to reduce itching.
As a topical analgesic, it is used to relieve minor aches and pains, such as muscle cramps, sprains, headaches and similar conditions, alone or combined with chemicals such as camphor, eucalyptus oil or capsaicin. In Europe, it tends to appear as a gel or a cream, while in the U.S., patches and body sleeves are very frequently used, e.g.: Tiger Balm, or IcyHot patches or knee/elbow sleeves.
As a penetration enhancer in transdermal drug delivery.
In decongestants for chest and sinuses (cream, patch or nose inhaler).
Examples: Vicks VapoRub, Mentholatum, Axe Brand, VapoRem, Mentisan.
In certain medications used to treat sunburns, as it provides a cooling sensation (then often associated with aloe).
Commonly used in oral hygiene products and bad-breath remedies, such as mouthwash, toothpaste, mouth and tongue sprays, and more generally as a food flavor agent; such as in chewing gum and candy.
In first aid products such as "mineral ice" to produce a cooling effect as a substitute for real ice in the absence of water or electricity (pouch, body patch/sleeve or cream).
In nonprescription products for short-term relief of minor sore throat and minor mouth or throat irritation e.g.: lip balms and cough medicines.
A recent study showed improvement in Alzheimer's symptoms and cognition improvements in mice.
Others
In aftershave products to relieve razor burn.
As a smoking tobacco additive in some cigarette brands, for flavor, and to reduce throat and sinus irritation caused by smoking. Menthol also increases nicotine receptor density, increasing the addictive potential of tobacco products.
As a pesticide against tracheal mites of honey bees.
In perfumery, menthol is used to prepare menthyl esters to emphasize floral notes (especially rose).
In various patches ranging from fever-reducing patches applied to children's foreheads to "foot patches" to relieve numerous ailments (the latter being much more frequent and elaborate in Asia, especially Japan: some varieties use "functional protrusions", or small bumps to massage one's feet as well as soothing them and cooling them down).
As an antispasmodic and smooth muscle relaxant in upper gastrointestinal endoscopy.
Organic chemistry
In organic chemistry, menthol is used as a chiral auxiliary in asymmetric synthesis. For example, sulfinate esters made from sulfinyl chlorides and menthol can be used to make enantiomerically pure sulfoxides by reaction with organolithium reagents or Grignard reagents. Menthol reacts with chiral carboxylic acids to give diastereomic menthyl esters, which are useful for chiral resolution.
It can be used as a catalyst for sodium production for the amateur chemist via the alcohol catalysed magnesium reduction process.
Menthol is potentially ergogenic (performance enhancing) for athletic performance in hot environments
Reactions
Menthol reacts in many ways like a normal secondary alcohol. It is oxidised to menthone by oxidising agents such as chromic acid, dichromate, or by calcium hypochlorite, in a green chemistry route. Under some conditions the oxidation using Cr(VI) compounds can go further and break open the ring. Menthol is easily dehydrated to give mainly 3-menthene, by the action of 2% sulfuric acid. Phosphorus pentachloride (PCl5) gives menthyl chloride.
History
In the West, menthol was first isolated in 1771, by the German, Hieronymus David Gaubius. Early characterizations were done by Oppenheim, Beckett, Moriya, and Atkinson. It was named by F. L. Alphons Oppenheim (1833–1877) in 1861.
Compendial status
United States Pharmacopeia 23
Japanese Pharmacopoeia 15
Food Chemicals Codex
Safety
The estimated lethal dose for menthol (and peppermint oil) in humans may be as low as LD=50–500 mg/kg. In the rat, 3300 mg/kg. In the mouse, 3400 mg/kg. In the cat, 800 mg/kg.
Survival after doses of 8 to 9 g has been reported. Overdose effects are abdominal pain, ataxia, atrial fibrillation, bradycardia, coma, dizziness, lethargy, nausea, skin rash, tremor, vomiting, and vertigo.
| Physical sciences | Terpenes and terpenoids | Chemistry |
317625 | https://en.wikipedia.org/wiki/Image%20scanner | Image scanner | An image scanner (often abbreviated to just scanner) is a device that optically scans images, printed text, handwriting, or an object and converts it to a digital image. The most common type of scanner used in the home and the office is the flatbed scanner, where the document is placed on a glass bed. A sheetfed scanner, which moves the page across an image sensor using a series of rollers, may be used to scan one page of a document at a time or multiple pages, as in an automatic document feeder. A handheld scanner is a portable version of an image scanner that can be used on any flat surface. Scans are typically downloaded to the computer that the scanner is connected to, although some scanners are able to store scans on standalone flash media (e.g., memory cards and USB drives).
Modern scanners typically use a charge-coupled device (CCD) or a contact image sensor (CIS) as the image sensor, whereas drum scanners, developed earlier and still used for the highest possible image quality, use a photomultiplier tube (PMT) as the image sensor. Document cameras, which use commodity or specialized high-resolution cameras, photograph documents all at once.
History
Precursors
Image scanners are considered the successors of early facsimile (fax) machines. The earliest attempt at a fax machine was patented in 1843 by the Scottish clockmaker Alexander Bain but never put into production. In his design, a metal stylus linked to a pendulum scans across a copper plate with a raised image. When the stylus makes contact with a raised part of the plate, it sends a pulse across a pair of wires to a receiver containing an electrode linked to another pendulum. A piece of paper impregnated with an electrochemically sensitive solution resides underneath the electrode and changes color whenever a pulse reaches the electrode. A gear advances the copper plate and paper in tandem with each swing of the pendulum; over time, the result is a perfect reproduction of the copper plate. In Bain's system, it is critical that the pendulums of the transceiver and receiver are in perfect step, or else the reproduced image will be distorted.
In 1847, the English physicist Frederick Bakewell developed the first working fax machine. Bakewell's machine was similar to Bain's but used a revolving drum coated in tinfoil, with non-conductive ink painted on the foil and a stylus that scans across the drum and sends a pulse down a pair of wires when it contacts a conductive point on the foil. The receiver contains an electrode that touches a sheet of chemically treated paper, which changes color when the electrode receives a pulse; the result is a reverse contrast (white-on-blue) reproduction of the original image. Bakewell's fax machine was marginally more successful than Bain's but suffered from the same synchronization issues. In 1862, Giovanni Caselli solved this with the pantelegraph, the first fax machine put into regular service. Largely based on Bain's design, it ensured complete synchronization by flanking the pendulums of both the transceiver and receiver between two magnetic regulators, which become magnetized with each swing of the pendulum and become demagnetized when the pendulum reaches the maxima and minima of each oscillation.
In 1893, the American engineer Elisha Gray introduced the telautograph, the first widely commercially successful fax machine that used linkage bars translating x- and y-axis motion at the receiver to scan a pen across the paper and strike it only when actuated by the stylus moving across the transceiver drum. Because it could use commodity stationery paper, it became popular in business and hospitals. In 1902, the German engineer Arthur Korn introduced the phototelautograph, a fax machine that used a light-sensitive selenium cell to scan a paper to be copied, instead of relying on a metallic drum and stylus. It was even more commercially successful than Gray's machine and became the basis for telephotography machines used by newspapers around the world from the early 1900s onward.
Analog era
Alexander Murray and Richard Morse invented and patented the first analog color scanner at Eastman Kodak in 1937. Intended for color separation at printing presses, their machine was an analog drum scanner that imaged a color transparency mounted in the drum, with a light source placed underneath the film, and three photocells with red, green, and blue color filters reading each spot on the transparency to translate the image into three electronic signals. In Murray and Morse's initial design, the drum was connected to three lathes that etched cyan, magenta, and yellow (CMY) halftone dots onto three offset cylinders directly. The rights to the patent were sold to Printing Developments Incorporated (P.D.I.) in 1946, who improved on the design by using a photomultiplier tube to image the points on the negative, which produced an amplified signal that was then fed to a single-purpose computer that processed the RGB signals into color-corrected cyan, magenta, yellow, and black (CMYK) values. The processed signals are then sent to four lathes that etch CMYK halftone dots onto the offset cylinders.
In 1948, Arthur Hardy of the Interchemical Corporation and F. L. Wurzburg of the Massachusetts Institute of Technology invented the first analog, color flatbed image scanner, intended for producing color-corrected lithographic plates from a color negative. In this system, three color-separated plates (of CMY values) are prepared from a color negative via dot etching and placed in the scanner bed. Above each plate are rigidly fixed, equidistant light beam projectors that focus a beam of light onto one corner of the plate. The entire bed with all three plates moves horizontally, back and forth, to reach the opposite corners of the plate; with each horiztonal oscillation of the bed, the bed moves down one step to cover the entire vertical area of the plate. While this is happening, the beam of light focused on a given spot on the plate gets reflected and bounced off to a photocell adjacent to the projector. Each photocell connects to an analog image processor, which evaluates the reflectance of the combined CMY values using Neugebauer equations and outputs a signal to a light projector hovering over a fourth, unexposed lithographic plate. This plate receives a color-corrected, continuous-tone dot-etch of either the cyan, magenta, or yellow values. The fourth plate is replaced with another unexposed plate, and the process repeats until three color-corrected plates, of cyan, magenta and yellow, are produced. In the 1950s, the Radio Corporation of America (RCA) took Hardy and Wurzburg's patent and replaced the projector-and-photocell arrangement with a video camera tube focusing on one spot of the plate.
Digital era
The first digital imaging system was the Bartlane system in 1920. Named after the pair who invented it, Harry G. Bartholomew and Maynard D. McFarlane, the Bartlane system used zinc plates etched with an image from a film negative projected at five different exposure levels to correspond to five quantization levels. All five plates are affixed to a long, motor-driven rotating cylinder, with five equidistant contacts scanning over each plate at the same starting position. The Bartlane system was initially used exclusively by telegraph, with the five-bit Baudot code used to transmit the grayscale digital image. In 1921, the system was modified for offline use, with a five-bit paper tape punch punching holes depending on whether its connections to the contacts are bridged or not. The result was a stored digital image with five gray levels. Reproduction of the image was achieved with a lamp passing over the punched holes, exposing five different intensities of light onto a film negative.
The first scanner to store its images digitally onto a computer was a drum scanner built in 1957 at the National Bureau of Standards (NBS, later NIST) by a team led by Russell A. Kirsch. It used a photomultiplier tube to detect light at a given point and produced an amplified signal that a computer could read and store into memory. The computer of choice at the time was the SEAC mainframe; the maximum horizontal resolution that the SEAC was capable of processing was 176 pixels. The first image ever scanned on this machine was a photograph of Kirsch's three-month-old son, Walden.
In 1969, Dacom introduced the 111 fax machine, which was the first digital fax machine to employ data compression using an on-board computer. It employed a flatbed design with a continuous feed capable of scanning up to letter paper in 1-bit monochrome (black and white).
The first flatbed scanner used for digital image processing was the Autokon 8400, introduced by ECRM Inc., a subsidiary of AM International, in 1975. The Autokon 8400 used a laser beam to scan pages up to 11 by 14 inches at a maximum resolution of 1000 lines per inch. Although it was only capable of scanning in 1-bit monochrome, the on-board processor was capable of halftoning, unsharp masking, contrast adjustment, and anamorphic distortions, among other features. The Autokon 8400 could either be connected to a film recorder to create a negative for producing plates or connected to a mainframe or minicomputer for further image processing and digital storage. The Autokon 8400 enjoyed widespread use in newspapers—ECRM shipped 1,000 units to newspaper publishers by 1985—but its limited resolution and maximum scan size made it unsuitable for commercial printing. In 1982, ECRM introduced the Autokon 8500, capable of scanning up to 1200 lines per inch. Four of ECRM's competitors introduced commercial flatbed scanners that year, including Scitex, Agfa-Gevaert, and Linotype-Hell, all of which were capable of scanning larger prints at higher resolutions. ECRM introduced the Autokon 1000DE in 1985 to address the shortcomings of the Autokon 8400/8500. The 1000DE (digital enhancement) used a microprocessor to produce the sharpening effect as against the 8400 which used analogue electronics and an optical method to create sharpening. The Autokon 1000DE had a touchpad rather than analogue rotary controls. The Autokon 1000DE had applications in both commercial and newspaper environments where only a single halftone was required ie black and white. Whilst typically the Autokon 8400 was a standalone output device that scanned and then output to either photosensitive, roll format bromide paper or film, the Autokon 1000DE was often connected to Apple Macintoshes or PCs via a dedicated interface such as those from HighWater Designs. The last Autokon was a wider format, online only device which utilised both a red and green laser to improve the response to the scanning of colour photographs.
In 1977, Raymond Kurzweil, of his start-up company Kurzweil Computer Products, released the Kurzweil Reading Machine, which was the first flatbed scanner with a charge-coupled device (CCD) imaging element. The Kurzweil Reading Machine was invented to assist blind people in reading books that had not been translated to braille. It comprised the image scanner and a Data General Nova minicomputer—the latter performing the image processing, optical character recognition (OCR), and speech synthesis.
The first scanners for personal computers appeared in the mid-1980s, starting with ThunderScan for the Macintosh in December 1984. Designed by Andy Hertzfeld and released by Thunderware Inc., the ThunderScan contains a specialized image sensor built into a plastic housing the same shape as the ink ribbon cartridge of Apple's ImageWriter printer. The ThunderScan slots into the ImageWriter's ribbon carrier and connects to both the ImageWriter and the Macintosh simultaneously. The ImageWriter's carriage, controlled by the ThunderScan, moves left-to-right to scan one 200-dpi (dots per inch) line at a time, with the carriage return serving to advance the scanner down the print to be scanned. The ThunderScan was the Macintosh's first scanner and sold well but operated very slowly and was only capable of scanning prints at 1-bit monochrome. In 1999, Canon iterated on this idea with the IS-22, a cartridge that fit into their inkjet printers to convert them into sheetfed scanners.
In early 1985, the first flatbed scanner for the IBM PC, the Datacopy Model 700, was released. Based on a CCD imaging element, the Model 700 was capable of scanning letter-sized documents at a maximum resolution of 200 dpi at 1-bit monochrome. The Model 700 came with a special interface card for connecting to the PC, and an optional, aftermarket OCR software card and software package were sold for the Model 700. In April 1985, LaserFAX Inc. introduced the first CCD-based color flatbed scanner, the SpectraSCAN 200, for the IBM PC. The SpectraSCAN 200 worked by placing color filters over the CCD and taking four passes (three for each primary color and one for black) per scan to build up a color reproduction. The SpectraSCAN 200 took between two and three minutes to produce a scan of a letter-sized print at 200-dpi; its grayscale counterpart, the DS-200, took only 30 seconds to make a scan at the same size and resolution.
The first relatively affordable flatbed scanner for personal computers appeared in February 1987 with Hewlett-Packard's ScanJet, which was capable of scanning 4-bit (64-shade) grayscale images at a maximum resolution of 300 dpi. By the beginning of 1988, the ScanJet had accounted for 27 percent of all scanner sales in terms of dollar volume, per Gartner Dataquest. In February 1989, the company introduced the ScanJet Plus, which increased the bit depth to 8 bits (256 shades) while costing only US$200 more than the original ScanJet's $1990 (). This led to a massive price drop in grayscale scanners with equivalent or lesser features in the market. The number of third-party developers producing software and hardware supporting these scanners jumped dramatically in turn, effectively popularizing the scanner for the personal computer user. By 1999, the cost of the average color-capable scanner had dropped to $300 (). That year, Computer Shopper declared 1999 "the year that scanners finally became a mainstream commodity".
Types
Flatbed
A flatbed scanner is a type of scanner that provides a glass bed (platen) on which the object to be scanned lies motionless. The scanning element moves vertically from under the glass, scanning either the entirety of the platen or a predetermined portion. The driver software for most flatbed scanners allows users to prescan their documents—in essence, to take a quick, low-resolution pass at a document in order to judge what area of the document should be scanned (if not the entirety of it), before scanning it at a higher resolution. Some flatbed scanners incorporate sheet-feeding mechanisms called automatic document feeders (ADFs) that use the same scanning element as the flatbed portion.
This type of scanner is sometimes called a reflective scanner, because it works by shining white light onto the object to be scanned and reading the intensity and color of light that is reflected from it, usually a line at a time. They are designed for scanning prints or other flat, opaque materials, but some have available transparency adapters, which—for a number of reasons—in most cases, are not very well suited to scanning film.
Sheetfed
A sheetfed scanner, also known as a document feeder, is a type of scanner that uses motor-driven rollers to move one single sheet of paper at a time past a stationary scanning element (two scanning elements, in the case of scanners with duplex functionality). Unlike flatbed scanners, sheetfed scanners are not equipped to scan bound material such as books or magazines, nor are they suitable for any material thicker than plain printer paper. Some sheetfed scanners, called automatic document feeders (ADFs), are capable of scanning several sheets in one session, although others only accept one page at a time. Some sheetfed scanners are portable, powered by batteries, and have their own storage, eventually transferring stored scans to a computer.
Handheld
A handheld scanner is a type of scanner that must be manually dragged or gilded by hand across the surface of the object to be scanned. Scanning documents in this manner requires a steady hand, as an uneven scanning rate produces distorted images. Some handheld scanners have an indicator light on the scanner for this purpose, actuating if the user is moving the scanner too fast. They typically have at least one button that starts the scan when pressed; it is held by the user for the duration of the scan. Some other handheld scanners have switches to set the optical resolution, as well as a roller, which generates a clock pulse for synchronization with the computer. Older hand scanners were monochrome, and produced light from an array of green LEDs to illuminate the image; later ones scan in monochrome or color, as desired. A hand scanner may also have a small window through which the document being scanned could be viewed. As hand scanners are much narrower than most normal document or book sizes, software (or the end user) needed to combine several narrow "strips" of scanned documents to produce the finished article.
Inexpensive, portable, battery-powered or USB-powered wand scanners and pen scanners, typically capable of scanning an area as wide as a normal letter and much longer, remain available . Some computer mice can also scan documents.
Drum
A drum scanner is a type of scanner that uses a clear, motor-driven rotating cylinder (drum) onto which a print, a film negative, a transparency, or any other flat object is taped or otherwise secured. A beam of light either projects past, or reflects off, the material to be scanned onto a series of mirrors, which focus the beam onto the drum scanner's photomultiplier tube (PMT). After one revolution, the beam of light moves down a single step. When scanning transparent media, such as negatives, a light beam is directed from within the cylinder onto the media; when scanning opaque items, a light beam from above is reflected off the surface of the media. When only one PMT is present, three passes of the image are required for a full-color RGB scan. When three PMTs are present, only a single pass is required.
The photomultiplier tubes of drum scanners offer superior dynamic range to that of CCD sensors. For this reason, drum scanners can extract more detail from very dark shadow areas of a transparency than flatbed scanners using CCD sensors. The smaller dynamic range of the CCD sensors (versus photomultiplier tubes) can lead to loss of shadow detail, especially when scanning very dense transparency film. Drum scanners are also able to resolve true detail in excess of 10000 dpi, producing higher-resolution scans than any CCD scanner.
Overhead
An overhead scanner is a type of scanner that places the scanning element in a housing on top of a vertical post, hovering above the document or object to be scanned, which lies stationary on an open-air bed. Chinon Industries patented a specific type of overhead scanner, which uses a rotating mirror to reflect the contents of the bed onto a linear CCD, in 1987. Although very flexible—allowing users to scan not only two-dimensional prints and documents but any 3D object, of any size—the Chinon design required the user to provide uniform illumination of the object to be scanned and was more cumbersome to set up.
A more modern type of overhead scanner is a document camera (also known as a video scanner), which uses a digital camera to capture a document all at once. Most document cameras output live video of the document and are usually reserved for displaying documents to a live audience, but they may also be used as replacements for image scanners, capturing a single frame of the output as an image file. Document cameras may even use the same APIs as scanners when connected to computers. A planetary scanner is a type of very-high-resolution document camera used for capturing certain fragile documents. A book scanner is another kind of document camera, pairing a digital camera with a scanning area defined by a mat to assist in scanning books. Some more advanced models of book scanners project a laser onto the page for calibration and software skew correction.
Film
A film scanner, also known as a slide scanner or a transparency scanner, is a type of specialized flatbed scanner specifically for scanning film negatives and slides. A typical film scanner works by passing a narrowly focused beam of light through the film and reading the intensity and color of the light that emerges. The lowest-cost dedicated film scanners can be had for less than $50, and they might be sufficient for modest needs. From there they inch up in staggered levels of quality and advanced features upward of five figures.
Portable
Image scanners are usually used in conjunction with a computer which controls the scanner and stores scans. Small portable scanners, either sheetfed or handheld and operated by batteries and with storage capability, are available for use away from a computer; stored scans can be transferred later. Many can scan both small documents such as business cards and till receipts, as well as letter-sized documents.
Software scanners
The higher-resolution cameras fitted to some smartphones can produce reasonable quality document scans by taking a photograph with the phone's camera and post-processing it with a scanning app, a range of which are available for most phone operating systems, to whiten the background of a page, correct perspective distortion so that the shape of a rectangular document is corrected, convert to black-and-white, etc. Many such apps can scan multiple-page documents with successive camera exposures and output them either as a single file or multiple-page files. Some smartphone scanning apps can save documents directly to online storage locations, such as Dropbox and Evernote, send via email, or fax documents via email-to-fax gateways.
Smartphone scanner apps can be broadly divided into three categories:
Document scanning apps primarily designed to handle documents and output PDF, and sometimes JPEG, files
Photo scanning apps that output JPEG files, and have editing functions useful for photo rather than document editing;
Barcode-like QR code scanning apps that then search the internet for information associated with the code.
Scanning elements
Charge-coupled device (CCD)
Scanners equipped with charge-coupled device (CCD) scanning elements require a sophisticated series of mirrors and lenses to reproduce an image, but the result of this complexity is a much higher-quality scan. Because CCDs have a much greater depth of field, they are more forgiving when it comes to scanning documents that are difficult to get perfectly flat against the platen (such as bound books).
Contact image sensor (CIS)
Scanners equipped with contact image sensor (CIS) scanning elements are designed to be in near-direct contact with the document to be scanned and thus do not require the complex optics of CCDs scanners. However, their depth of field is much worse, resulting in blurry scans if the scanned document is not perfectly flush against the platten. Because the sensors require far less power than CCD scanners, CIS scanners are able to be manufactured down to a low cost and are typically much lighter in weight and depth than CCD scanners.
Photomultiplier tube (PMT)
Scanners equipped with photomultiplier tubes (PMT) are nearly exclusively drum scanners.
Scan quality
Color scanners typically read RGB (red-green-blue) color data from the array. This data is then processed with some proprietary algorithm to correct for different exposure conditions, and sent to the computer via the device's input/output interface (usually USB, previous to which was SCSI or bidirectional parallel port in older units).
Color depth varies depending on the scanning array characteristics, but is usually at least 24 bits. High-quality models have 36-48 bits of color depth.
Another qualifying parameter for a scanner is its resolution, measured in pixels per inch (ppi), sometimes more accurately referred to as samples per inch (spi). Instead of using the scanner's true optical resolution, the only meaningful parameter, manufacturers like to refer to the interpolated resolution, which is much higher thanks to software interpolation. , a high-end flatbed scanner can scan up to 5400 ppi and drum scanners have an optical resolution of between 3000 and 24000 ppi.
Effective resolution refers to the true resolution of a scanner, and is determined by using a resolution test chart. The effective resolution of most all consumer flatbed scanners is considerably lower than the manufactures' given optical resolution.
Manufacturers often claim interpolated resolutions as high as 19200 ppi; but such numbers carry little meaningful value because the number of possible interpolated pixels is unlimited, and doing so does not increase the level of captured detail.
The size of the file created increases with the square of the resolution; doubling the resolution quadruples the file size. A resolution must be chosen that is within the capabilities of the equipment, preserves sufficient detail, and does not produce a file of excessive size. The file size can be reduced for a given resolution by using "lossy" compression methods such as JPEG, at some cost in quality. If the best possible quality is required lossless compression should be used; reduced-quality files of smaller size can be produced from such an image when required (e.g., image designed to be printed on a full page, and a much smaller file to be displayed as part of a fast-loading web page).
Purity can be diminished by scanner noise, optical flare, poor analog to digital conversion, scratches, dust, Newton's rings, out-of-focus sensors, improper scanner operation, and poor software. Drum scanners are said to produce the purest digital representations of the film, followed by high-end film scanners that use the larger Kodak Tri-Linear sensors.
The third important parameter for a scanner is its dynamic range (also known as density range). A high-density range means that the scanner is able to record shadow details and brightness details in one scan. Density of film is measured on a base 10 log scale and varies between 0.0 (transparent) and 5.0, about 16 stops. Density range is the space taken up in the 0 to 5 scale, and Dmin and Dmax denote where the least dense and most dense measurements on a negative or positive film. The density range of negative film is up to 3.6d, while slide film dynamic range is 2.4d. Color negative density range after processing is 2.0d thanks to the compression of the 12 stops into a small density range. Dmax will be the densest on slide film for shadows, and densest on negative film for highlights. Some slide films can have a Dmax close to 4.0d with proper exposure, and so can black-and-white negative film.
Consumer-level flatbed photo scanners have a dynamic range in the 2.0–3.0 range, which can be inadequate for scanning all types of photographic film, as Dmax can be and often is between 3.0d and 4.0d with traditional black-and-white film. Color film compresses its 12 stops of a possible 16 stops (film latitude) into just 2.0d of space via the process of dye coupling and removal of all silver from the emulsion. Kodak Vision 3 has 18 stops. So, color-negative film scans the easiest of all film types on the widest range of scanners. Because traditional black-and-white film retains the image creating silver after processing, density range can be almost twice that of color film. This makes scanning traditional black-and-white film more difficult and requires a scanner with at least a 3.6d dynamic range, but also a Dmax between 4.0d to 5.0d. High-end (photo lab) flatbed scanners can reach a dynamic range of 3.7, and Dmax around 4.0d. Dedicated film scanners have a dynamic range between 3.0d–4.0d. Office document scanners can have a dynamic range of less than 2.0d. Drum scanners have a dynamic range of 3.6–4.5.
For scanning film, is a technique used to remove the effects of dust and scratches on images scanned from film; many modern scanners incorporate this feature. It works by scanning the film with infrared light; the dyes in typical color film emulsions are transparent to infrared light, but dust and scratches are not, and block infrared; scanner software can use the visible and infrared information to detect scratches and process the image to greatly reduce their visibility, considering their position, size, shape, and surroundings. Scanner manufacturers usually have their own names attached to this technique. For example, Epson, Minolta, Nikon, Konica Minolta, Microtek, and others use Digital ICE, while Canon uses its own system, FARE (Film Automatic Retouching and Enhancement). Plustek uses LaserSoft Imaging iSRD. Some independent software developers design infrared cleaning tools.
By combining full-color imagery with 3D models, modern hand-held scanners are able to completely reproduce objects electronically. The addition of 3D color printers enables accurate miniaturization of these objects, with applications across many industries and professions.
For scanner apps, the scan quality is highly dependent on the quality of the phone camera and on the framing chosen by the user of the app.
Connectivity
Scans must virtually always be transferred from the scanner to a computer or information storage system for further processing or storage. There are two basic issues: (1) how the scanner is physically connected to the computer and (2) how the application retrieves the information from the scanner.
Direct connection
The file size of a scan can go up to about 100 MB for a 600 dpi, 23 × 28 cm (slightly larger than A4 paper) uncompressed 24-bit image. Scanned files must be transferred and stored. Scanners can generate this volume of data in a matter of seconds, making a fast connection desirable.
Scanners communicate to their host computer using one of the following physical interfaces, listing roughly from slow to fast:
Parallel port – Connecting through a parallel port is the slowest common transfer method. Early scanners had parallel port connections that could not transfer data faster than 70 kilobytes/second. The primary advantage of the parallel port connection was economic and user skill level: it avoided adding an interface card to the computer.
GPIB – General Purpose Interface Bus. Certain drum scanners like the Howtek D4000 featured both a SCSI and GPIB interface. The latter conforms to the IEEE-488 standard, introduced in the mid-1970s. The GPIB interface has only been used by a few scanner manufacturers, mostly serving the DOS/Windows environment. For Apple Macintosh systems, National Instruments provided a NuBus GPIB interface card.
Small Computer System Interface (SCSI) – SCSI is rarely used since the early 21st century, supported only by computers with a SCSI interface, either on a card or built-in. During the evolution of the SCSI standard, speeds increased. Widely available and easily set up USB and Firewire largely supplanted SCSI.
Universal Serial Bus (USB) – USB scanners can transfer data quickly. The early USB 1.1 standard could transfer data at 1.5 megabytes per second (slower than SCSI), but the later USB 2.0/3.0 standards can transfer at more than 20/60 megabytes per second in practice.
FireWire – Also known as IEEE-1394, FireWire is an interface of comparable speed to USB 2.0. Possible FireWire speeds are 25, 50, and 100, 400, and 800 megabits per second, but devices may not support all speeds.
Proprietary interfaces – Bespoke interfaces were used on some early scanners that used a proprietary interface card rather than a standard interface.
Indirect connection
During the early 1990s professional flatbed scanners were available over a local computer network. This proved useful to publishers, print shops, etc. This functionality largely fell out of use as the cost of flatbed scanners reduced enough to make sharing unnecessary.
From 2000 all-in-one multi-purpose devices became available which were suitable for both small offices and consumers, with printing, scanning, copying, and fax capability in a single apparatus that can be made available to all members of a workgroup.
Battery-powered portable scanners store scans on internal memory; they can later be transferred to a computer either by direct connection, typically USB, or in some cases a memory card may be removed from the scanner and plugged into the computer.
Applications programming interface
A raster image editor must be able to communicate with a scanner. There are many different scanners, and many of those scanners use different protocols. In order to simplify applications programming, some application programming interfaces (APIs) were developed. The API presents a uniform interface to the scanner. This means that the application does not need to know the specific details of the scanner in order to access it directly. For example, Adobe Photoshop supports the TWAIN standard; therefore in theory Photoshop can acquire an image from any scanner that has a TWAIN driver.
In practice, there are often problems with an application communicating with a scanner. Either the application or the scanner manufacturer (or both) may have faults in their implementation of the API.
Typically, the API is implemented as a dynamically linked library. Each scanner manufacturer provides software that translates the API procedure calls into primitive commands that are issued to a hardware controller (such as the SCSI, USB, or FireWire controller). The manufacturer's part of the API is commonly called a device driver, but that designation is not strictly accurate: the API does not run in kernel mode and does not directly access the device. Rather the scanner API library translates application requests into hardware requests.
Common scanner software API include:
TWAIN – An API used by most scanners. Originally used for low-end and home-use equipment, it is now widely used for large-volume scanning.
SANE (Scanner Access Now Easy) – A free/open-source API for accessing scanners. Originally developed for Unix and Linux operating systems, it has been ported to OS/2, Mac OS X, and Microsoft Windows. Unlike TWAIN, SANE does not handle the user interface. This allows batch scans and transparent network access without any special support from the device driver.
Windows Image Acquisition (WIA) – An API provided by Microsoft for use on Microsoft Windows.
Image and Scanner Interface Specification (ISIS) – Created by Pixel Translations, which still uses SCSI-2 for performance reasons, ISIS is used by large, departmental-scale, machines.
Bundled applications
Although no software beyond a scanning utility is a feature of any scanner, many scanners come bundled with software. Typically, in addition to the scanning utility, some type of raster image editor (such as Photoshop or GIMP) and optical character recognition (OCR) software are supplied. OCR software converts graphical images of text into standard text that can be edited using common word-processing and text-editing software; accuracy is rarely perfect.
Output data
Some scanners, especially those designed for scanning printed documents, only work in black and white, but most modern scanners work in color. For the latter, the scanned result is a non-compressed RGB image, which can be transferred to a computer's memory. The color output of different scanners is not the same due to the spectral response of their sensing elements, the nature of their light source, and the correction applied by the scanning software. While most image sensors have a linear response, the output values are usually gamma-compressed. Some scanners compress and clean up the image using embedded firmware. Once on the computer, the image can be processed with a raster graphics editor (such as Photoshop) and saved on a storage device (such as a hard disk).
Scans may be stored uncompressed in image file formats such as BMP; stored losslessly compressed in file formats such as TIFF and PNG; stored lossy-compressed in file formats such as JPEG; or stored as embedded images or converted to vector graphics within a PDF. Optical character recognition (OCR) software allows a scanned image of text to be converted into editable text with reasonable accuracy, so long as the text is cleanly printed and in a typeface and size that can be read by the software. OCR capability may be integrated into the scanning software, or the scanned image file can be processed with a separate OCR program.
Specific uses
Document processing
Document processing requirements differ from those of image scanning. These requirements include scanning speed, automated paper feed, and the ability to automatically scan both the front and the back of a document. On the other hand, image scanning typically requires the ability to handle fragile and or three-dimensional objects as well as scan at a much higher resolution.
Document scanners have document feeders, usually larger than those sometimes found on copiers or all-purpose scanners. Scans are made at high speed, from 20 up to 420 pages per minute, often in grayscale, although many scanners support color. Many scanners can scan both sides of double-sided originals (duplex operation). Sophisticated document scanners have firmware or software that cleans up scans of text as they are produced, eliminating accidental marks and sharpening type; this would be unacceptable for photographic work, where marks cannot reliably be distinguished from desired fine detail. Files created are compressed as they are made.
The resolution used is usually from 150 to 300 dpi, although the hardware may be capable of 600 or higher resolution; this produces images of text good enough to read and for OCR, without the higher demands on storage space required by higher-resolution images.
Document scans are often processed using OCR technology to create editable and searchable files. Most scanners use ISIS or TWAIN device drivers to scan documents into TIFF format so that the scanned pages can be fed into a document management system that will handle the archiving and retrieval of the scanned pages. Lossy JPEG compression, which is very efficient for pictures, is undesirable for text documents, as slanted straight edges take on a jagged appearance, and solid black (or other color) text on a light background compresses well with lossless compression formats.
While paper feeding and scanning can be done automatically and quickly, preparation and indexing are necessary and require much work by humans. Preparation involves manually inspecting the papers to be scanned and making sure that they are in order, unfolded, without staples or anything else that might jam the scanner. Additionally, some industries such as legal and medical may require documents to have Bates Numbering or some other mark giving a document identification number and date/time of the document scan.
Indexing involves associating relevant keywords to files so that they can be retrieved by content. This process can sometimes be automated to some extent, but it often requires manual labour performed by data-entry clerks. One common practice is the use of barcode-recognition technology: during preparation, barcode sheets with folder names or index information are inserted into the document files, folders, and document groups. Using automatic batch scanning, the documents are saved into appropriate folders, and an index is created for integration into document management systems.
A specialized form of document scanning is book scanning. Technical difficulties arise from the books usually being bound and sometimes fragile and irreplaceable, but some manufacturers have developed specialized machinery to deal with this. Often special robotic mechanisms are used to automate the page-turning and scanning process.
Other uses
Flatbed scanners have been used as digital backs for large-format cameras to create high-resolution digital images of static subjects. A modified flatbed scanner has been used for documentation and quantification of thin layer chromatograms detected by fluorescence quenching on silica gel layers containing an ultraviolet (UV) indicator. The ChromImage is allegedly the first commercial flatbed scanner densitometer. It enables acquisition of TLC plate images and quantification of chromatograms by use of Galaxie-TLC software. Other than being turned into densitometers, flatbed scanners were also turned into colorimeters using different methods. Trichromatic Color Analyser is allegedly the first distributable system using a flatbed scanner as a tristimulus colorimetric device.
Flatbed scanners may also be used to create artwork directly, in a practice known as scanography.
In the biomedical research field, detection devices for DNA microarrays are also referred to as scanners. These scanners are high-resolution systems (up to 1 μm/pixel), similar to microscopes. Detection is performed using CCDs or photomultiplier tubes.
In pathology, scanners are used to capture glass slides with tissue from biopsies and other kinds of sampling, allowing for various methods of digital pathology such as telepathology and the application of artificial intelligence for interpretation.
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317695 | https://en.wikipedia.org/wiki/Moulting | Moulting | In biology, moulting (British English), or molting (American English), also known as sloughing, shedding, or in many invertebrates, ecdysis, is a process by which an animal casts off parts of its body to serve some beneficial purpose, either at specific times of the year, or at specific points in its life cycle.
In medieval times, it was also known as "mewing" (from the French verb "muer", to moult), a term that lives on in the name of Britain's Royal Mews where the King's hawks used to be kept during moulting time before becoming horse stables after Tudor times.
Moulting can involve shedding the epidermis (skin), pelage (hair, feathers, fur, wool), or other external layer. In some groups, other body parts may be shed, for example, the entire exoskeleton in arthropods, including the wings in some insects.
Examples
In birds
In birds, moulting is the periodic replacement of feathers by shedding old feathers while producing new ones. Feathers are dead structures at maturity which are gradually abraded and need to be replaced. Adult birds moult at least once a year, although many moult twice and a few three times each year. It is generally a slow process: birds rarely shed all their feathers at any one time. The bird must retain sufficient feathers to regulate its body temperature and repel moisture. The number and area of feathers that are shed varies. In some moulting periods, a bird may renew only the feathers on the head and body, shedding the wing and tail feathers during a later moulting period.
Some species of bird become flightless during an annual "wing moult" and must seek a protected habitat with a reliable food supply during that time. While the plumage may appear thin or uneven during the moult, the bird's general shape is maintained despite the loss of apparently many feathers; bald spots are typically signs of unrelated illnesses, such as gross injuries, parasites, feather pecking (especially in commercial poultry), or (in pet birds) feather plucking. Some birds will drop feathers, especially tail feathers, in what is called a "fright moult".
The process of moulting in birds is as follows: First, the bird begins to shed some old feathers, then pin feathers grow in to replace the old feathers. As the pin feathers become full feathers, other feathers are shed. This is a cyclical process that occurs in many phases. It is usually symmetrical, with feather loss equal on each side of the body. Because feathers make up 4–12% of a bird's body weight, it takes a large amount of energy to replace them.
For this reason, moults often occur immediately after the breeding season, but while food is still abundant. The plumage produced during this time is called postnuptial plumage. Prenuptial moulting occurs in red-collared widowbirds where the males replace their nonbreeding plumage with breeding plumage. It is thought that large birds can advance the moult of severely damaged feathers.
Determining the process birds go through during moult can be useful in understanding breeding, migration and foraging strategies. One non-invasive method of studying moult in birds is through using field photography. The evolutionary and ecological forces driving moult can also be investigated using intrinsic markers such as stable hydrogen isotope (δ2H) analysis. In some tropical birds, such as the common bulbul, breeding seasonality is weak at the population level, instead moult can show high seasonality with individuals probably under strong selection to match moult with peak environmental conditions.
A 2023 paleontological analysis concluded that moulting probably evolved late in the evolutionary lineage of birds.
Forced moulting
In some countries, flocks of commercial layer hens are force-moulted to reinvigorate egg-laying. This usually involves complete withdrawal of their food and sometimes water for 7–14 days or up to 28 days under experimental conditions, which presumably reflect standard farming practice in some countries. This causes a body weight loss of 25 to 35%, which stimulates the hen to lose her feathers, but also reinvigorates egg-production.
Some flocks may be force-moulted several times. In 2003, more than 75% of all flocks were force-moulted in the US. Other methods of inducing a moult include low-density diets (e.g. grape pomace, cotton seed meal, alfalfa meal) or dietary manipulation to create an imbalance of a particular nutrient(s). The most important among these include manipulation of minerals including sodium (Na), calcium (Ca), iodine (I) and zinc (Zn), with full or partially reduced dietary intakes.
In reptiles and amphibians
Squamates periodically engage in moulting, as their skin is scaly. The most familiar example of moulting in such reptiles is when snakes "shed their skin". This is usually achieved by the snake rubbing its head against a hard object, such as a rock (or between two rocks) or piece of wood, causing the already stretched skin to split.
At this point, the snake continues to rub its skin on objects, causing the end nearest the head to peel back on itself, until the snake is able to crawl out of its skin, effectively turning the moulted skin inside-out. This is similar to how one might remove a sock from one's foot by grabbing the open end and pulling it over itself. The snake's skin is often left in one piece after the moulting process, including the discarded brille (ocular scale), so that the moult is vital for maintaining the animal's quality of vision. The skins of lizards, in contrast, generally fall off in pieces.
Both frogs and salamanders moult regularly and consume the skin, with some species moulting in pieces and others in one piece.
In arthropods
In arthropods, such as insects, arachnids and crustaceans, moulting is the shedding of the exoskeleton, which is often called its shell, typically to let the organism grow. This process is called ecdysis. Most Arthropoda with soft, flexible skins also undergo ecdysis. Ecdysis permits metamorphosis, the sometimes radical difference between the morphology of successive instars.
A new skin can replace structures, such as by providing new external lenses for eyes. The new exoskeleton is initially soft but hardens after the moulting of the old exoskeleton. The old exoskeleton is called an exuviae. While moulting, insects cannot breathe. In the crustacean Ovalipes catharus molting must occur before they mate.
In dogs
Most dogs moult twice each year, in the spring and autumn, depending on the breed, environment and temperature. Dogs shedding much more than usual are known as "blow coats" or "blowing coats".
Gallery
| Biology and health sciences | Animal ontogeny | null |
317900 | https://en.wikipedia.org/wiki/Clothes%20dryer | Clothes dryer | A clothes dryer (tumble dryer, drying machine, or simply dryer) is a powered household appliance that is used to remove moisture from a load of clothing, bedding and other textiles, usually after they are washed in the washing machine.
Many dryers consist of a rotating drum called a "tumbler" through which heated air is circulated to evaporate moisture while the tumbler is rotated to maintain air space between the articles. Using such a machine may cause clothes to shrink or become less soft (due to loss of short soft fibers). A simpler non-rotating machine called a "drying cabinet" may be used for delicate fabrics and other items not suitable for a tumble dryer. Other machines include steam to de-shrink clothes and avoid ironing.
Tumble dryers
Tumble dryers continuously draw in the ambient air around them and heat it before passing it through the tumbler. The resulting hot, humid air is usually vented outside to make room for more air to continue the drying process.
Tumble dryers are sometimes integrated with a washing machine, in the form of washer-dryer combos, which are essentially a front loading washing machine with an integrated dryer or (in the US) a laundry center, which stacks the dryer on top of the washer and integrates the controls for both machines into a single control panel. Often the washer and dryer functions will have a different capacity, with the dryer usually having a lower capacity than the washer. Tumble dryers can also be top loading, in which the drum is loaded from the top of the machine and the drum's end supports are in the left and right sides, instead of the more conventional front and rear. They can be as thin as in width, and may include detachable stationary racks for drying items like plush toys and footwear.
Ventless dryers
Spin dryers
These centrifuge machines simply spin their drums much faster than a typical washer could, in order to extract additional water from the load. They may remove more water in two minutes than a heated tumbler dryer can in twenty, thus saving significant amounts of time and energy. Although spinning alone will not completely dry clothing, this additional step saves a worthwhile amount of time and energy for large laundry operations such as those of hospitals.
Condenser dryers
Just as in a tumble dryer, condenser or condensation dryers pass heated air through the load. However, instead of exhausting this air, the dryer uses a heat exchanger to cool the air and condense the water vapor into either a drain pipe or a collection tank. The drier air is run through the loop again. The heat exchanger typically uses ambient air as its coolant, therefore the heat produced by the dryer will go into the immediate surroundings instead of the outside, increasing the room temperature. In some designs, cold water is used in the heat exchanger, eliminating this heating, but requiring increased water usage.
In terms of energy use, condenser dryers typically require around 2 kilowatt hours (kW⋅h) of energy per average load.
Because the heat exchange process simply cools the internal air using ambient air (or cold water in some cases), it will not dry the air in the internal loop to as low a level of humidity as typical fresh, ambient air. As a consequence of the increased humidity of the air used to dry the load, this type of dryer requires somewhat more time than a tumble dryer. Condenser dryers are a particularly attractive option where long, intricate ducting would be required to vent the dryer.
Heat pump dryers
A closed-cycle heat pump clothes dryer uses a heat pump to dehumidify the processing air. Such dryers typically use under half the energy per load of a condenser dryer.
Whereas condensation dryers use a passive heat exchanger cooled by ambient air, these dryers use a heat pump. The hot, humid air from the tumbler is passed through a heat pump where the cold side condenses the water vapor into either a drain pipe or a collection tank and the hot side reheats the air afterward for re-use. In this way not only does the dryer avoid the need for ducting, but it also conserves much of its heat within the dryer instead of exhausting it into the surroundings. Heat pump dryers can, therefore, use up to 50% less energy required by either condensation or conventional electric dryers. Heat pump dryers use about 1 kW⋅h of energy to dry an average load instead of 2 kW⋅h for a condenser dryer, or from 3 to 9 kW⋅h, for a conventional electric dryer. Domestic heat pump dryers are designed to work in typical ambient temperatures from . Below , drying times significantly increase.
As with condensation dryers, the heat exchanger will not dry the internal air to as low a level of humidity as the typical ambient air. With respect to ambient air, the higher humidity of the air used to dry the clothes has the effect of increasing drying times; however, because heat pump dryers conserve much of the heat of the air they use, the already-hot air can be cycled more quickly, possibly leading to shorter drying times than tumble dryers, depending on the model.
Mechanical steam compression dryers
A new type of dryer in development, these machines are a more advanced version of heat pump dryers. Instead of using hot air to dry the clothing, mechanical steam compression dryers use water recovered from the clothing in the form of steam. First, the tumbler and its contents are heated to . The wet steam that results purges the system of air and is the only remaining atmosphere in the tumbler.
As wet steam exits the tumbler, it is mechanically compressed (hence the name) to extract water vapor and transfer the heat of vaporization to the remaining gaseous steam. This pressurized, gaseous steam is then allowed to expand, and is superheated before being injected back into the tumbler where its heat causes more water to vaporize from the clothing, creating more wet steam and restarting the cycle.
Like heat pump dryers, mechanical steam compression dryers recycle much of the heat used to dry the clothes, and they operate in a very similar range of efficiency as heat pump dryers. Both types can be over twice as efficient as conventional tumble dryers. The considerably higher temperatures used in mechanical steam compression dryers result in drying times on the order of half as long as those of heat pump dryers.
Convectant drying
Marketed by some manufacturers as a "static clothes drying technique", convectant dryers simply consist of a heating unit at the bottom, a vertical chamber, and a vent at top. The unit heats air at the bottom, reducing its relative humidity, and the natural tendency of hot air to rise brings this low-humidity air into contact with the clothes. This design is slower than conventional tumble dryers, but relatively energy-efficient if well-implemented. It works particularly well in cold and humid environments, where it dries clothes substantially faster than line-drying. In hot and dry weather, the performance delta over line-drying is negligible.
Given that this is a relatively simple and cheap technique to materialize, most consumer products showcase the added benefit of portability and/or modularity. Newer designs implement a fan heater at the bottom to pump hot air into the vertical drying rack chamber. Temperatures in excess of can be reached inside these "hot air balloons," yet lint, static cling, and shrinkage are minimal. Upfront cost is significantly lower than tumble, condenser and heat pump designs.
If used in combination with washing machines featuring fast spin cycles (800+ rpm) or spin dryers, the cost-effectiveness of this technique has the potential to render tumble dryer-like designs obsolete in single-person and small family households. One disadvantage is that the moisture from the clothes is released into the immediate surroundings. Proper ventilation or a complementary dehumidifier is recommended for indoor use. It also cannot compete with the tumble dryer's capacity to dry multiple loads of wet clothing in a single day.
Solar clothes dryer
The solar dryer is a box-shaped stationary construction which encloses a second compartment where the clothes are held. It uses the sun's heat without direct sunlight reaching the clothes. Alternatively, a solar heating box may be used to heat air that is driven through a conventional tumbler dryer.
Microwave dryers
Japanese manufacturers have developed highly efficient clothes dryers that use microwave radiation to dry the clothes (though a vast majority of Japanese air dry their laundry). Most of the drying is done using microwaves to evaporate the water, but the final drying is done by convection heating, to avoid problems of arcing with metal pieces in the laundry. There are a number of advantages: shorter drying times (25% less), energy savings (17–25% less), and lower drying temperatures. Some analysts think that the arcing and fabric damage is a factor preventing microwave dryers from being developed for the US market.
Ultrasonic dryers
Ultrasonic dryers use high-frequency signals to drive piezoelectric actuators in order to mechanically shake the clothes, releasing water in the form of a mist which is then removed from the drum. They have the potential to significantly cut energy consumption while needing only one-third of the time needed by a conventional electric dryer for a given load. They also do not have the same issues related with lint in most other types of dryers.
Hybrid dryers
Some manufacturers, like LG Electronics and Whirlpool, have introduced hybrid dryers, that offer the user the option of using either a heat pump or a traditional electric heating element for drying the user's clothes. Hybrid dryers can also use a heat pump and a heating element at the same time to dry clothes faster.
Static electricity
Clothes dryers can cause static cling through the triboelectric effect. This can be a minor nuisance and is often a symptom of over-drying textiles to below their equilibrium moisture level, particularly when using synthetic materials. Fabric conditioning products such as dryer sheets are marketed to dissipate this static charge, depositing surfactants onto the fabric load by mechanical abrasion during tumbling. Modern dryers often have improved temperature and humidity sensors and electronic controls which aim to stop the drying cycle once textiles are sufficiently dry, avoiding over-drying and the static charge and energy wastage this causes.
Pest control use
Drying at a minimum of heat for thirty minutes kills many parasites including house dust mites, bed bugs, and scabies mites and their eggs; a bit more than ten minutes kills ticks. Simply washing drowns dust mites, and exposure to direct sunlight for three hours kills their eggs.
Lint build-up (tumble dryers)
Moisture and lint are byproducts of the tumble drying process and are pulled from the drum by a fan motor and then pushed through the remaining exhaust conduit to the exterior termination fitting. Typical exhaust conduit comprises flex transition hose found immediately behind the dryer, the rigid galvanized pipe and elbow fittings found within the wall framing, and the vent duct hood found outside the house.
A clean, unobstructed dryer vent improves both the efficiency and safety of the dryer. As the dryer duct pipe becomes partially obstructed and filled with lint, drying time markedly increases and causes the dryer to waste energy. A blocked vent increases the internal temperature and may result in a fire. Clothes dryers are one of the more costly home appliances to operate.
Several factors can contribute to or accelerate rapid lint build-up. These include long or restrictive ducts, bird or rodent nests in the termination, crushed or kinked flex transition hose, terminations with screen-like features, and condensation within the duct due to un-insulated ducts traveling through cold spaces such as a crawl space or attic. If plastic flaps are at the outside end of the duct, one may be able to flex, bend, and temporarily remove the plastic flaps, clean the inside surface of the flaps, clean the last foot or so of the duct, and reattach the plastic flaps. The plastic flaps keep insects, birds, and snakes out of the dryer vent pipe. During cold weather, the warm wet air condenses on the plastic flaps, and minor trace amounts of lint sticks to the wet inside part of the plastic flaps at the outside of the building.
Ventless dryers include multi-stage lint filtration systems and some even include automatic evaporator and condenser cleaning functions that can run even while the dryer is running. The evaporator and condenser are usually cleaned with running water. These systems are necessary, in order to prevent lint from building up inside the dryer and evaporator and condenser coils.
Aftermarket add-on lint and moisture traps can be attached to the dryer duct pipe, on machines originally manufactured as outside-venting, to facilitate installation where an outside vent is not available. Increased humidity at the location of installation is a drawback to this method.
Safety
Dryers expose flammable materials to heat. Underwriters Laboratories recommends cleaning the lint filter after every cycle for safety and energy efficiency, provision of adequate ventilation, and cleaning of the duct at regular intervals. UL also recommends that dryers not be used for glass fiber, rubber, foam or plastic items, or any item that has had a flammable substance spilled on it.
In the United States, an estimate from the US Fire Administration in a 2012 report estimated that from 2008 to 2010, fire departments responded to an estimated 2,900 clothes dryer fires in residential buildings each year across the nation. These fires resulted in an annual average loss of 5 deaths, 100 injuries, and $35 million in property loss. The Fire Administration attributes "Failure to clean" (34%) as the leading factor contributing to clothes dryer fires in residential buildings, and observed that new home construction trends place clothes dryers and washing machines in more hazardous locations away from outside walls, such as in bedrooms, second-floor hallways, bathrooms, and kitchens.
To address the problem of clothes dryer fires, a fire suppression system can be used with sensors to detect the change in temperature when a blaze starts in a dryer drum. These sensors then activate a water vapor mechanism to put out the fire.
Environmental impact
The environmental impact of clothes dryers is especially severe in the US and Canada, where over 80% of all homes have a clothes dryer. According to the US Environmental Protection Agency, if all residential clothes dryers sold in the US were energy efficient, "the utility cost savings would grow to more than $1.5 billion each year and more than 10 billion kilograms (22 billion pounds) of annual greenhouse gas emissions would be prevented”.
Clothes dryers are second only to refrigerators and freezers as the largest residential electrical energy consumers in America.
In the European Union, the EU energy labeling system is applied to dryers; dryers are classified with a label from A+++ (best) to G (worst) according to the amount of energy used per kilogram of clothes (kW⋅h/kg). Sensor dryers can automatically sense that clothes are dry and switch off. This means over-drying is not as frequent. Most of the European market sells sensor dryers now, and they are normally available in condenser and vented dryers.
History
A hand-cranked clothes dryer was created in 1800 by M. Pochon from France. Henry W. Altorfer invented and patented an electric clothes dryer in 1937. J. Ross Moore, an inventor from North Dakota, developed designs for automatic clothes dryers and published his design for an electrically operated dryer in 1938. Industrial designer Brooks Stevens developed an electric dryer with a glass window in the early 1940s.
| Technology | Household appliances | null |
318020 | https://en.wikipedia.org/wiki/Megamouth%20shark | Megamouth shark | The megamouth shark (Megachasma pelagios) is a species of deepwater shark. Rarely seen by humans, it measures around long and is the smallest of the three extant filter-feeding sharks alongside the relatively larger whale shark and basking shark. According to Sharkman's World Organization a total of 286 specimens have been observed or caught since its discovery in 1976. Like the other two planktivorous sharks, it swims with its mouth wide open, filtering water for plankton and jellyfish. It is recognizable from its large head with rubbery lips. The megamouth is so unlike any other type of shark that it is usually considered to be the sole extant species in the family Megachasmidae, though some scientists have suggested it may belong in the family Cetorhinidae.
Taxonomy and evolution
The first megamouth shark was captured on November 15, 1976, about 25 miles northeast of Kahuku, Hawaii, when it became entangled in the sea anchor of United States Navy ship AFB-14 at a depth of about 165 m (541 ft). The species was identified as being of a new genus within the planktivorous shark species. Examination of the , specimen by Leighton Taylor showed it to be an entirely unknown type of shark, making it – along with the coelacanth – one of the more sensational discoveries in 20th-century ichthyology. The pectoral fin of the megamouth shark was studied, along with its skeletal and muscular system, to show its phylogenetic relationship to the other two sharks.
, only 99 megamouth specimens had been caught or sighted. They have been found in the Pacific, Atlantic, and Indian Oceans. Japan, the Philippines and Taiwan have each yielded at least 10 specimens, the most of any single area, amounting to more than half the worldwide total. Specimens have also been sighted in or come out of the waters near Hawaii, California, Mexico, Indonesia, Australia, Brazil, Senegal, South Africa, Puerto Rico, Ecuador, and possibly Vietnam.
Researchers have predicted the feeding patterns of megamouth sharks in relation to the other two planktivorous sharks; the three planktivorous sharks have ram feeding in common, as it evolved from ram feeding swimming-type ancestors that developed their filtering mechanism to capture small prey like plankton. In addition to the living M. pelagios, however, two extinct megamouth species – the Priabonian M. alisonae and the Oligocene–Miocene M. applegatei – have also recently been proposed on the basis of fossilized tooth remains. An early ancestor of the recent species Megachasma pelagios was reported from the early Miocene (Burdigalian) of Belgium. However, the Cretaceous-aged M. comanchensis has been recently reclassified as an odontaspid shark in the genus Pseudomegachasma, and is in fact unrelated to the megamouth shark despite similar teeth morphology. The megamouth's filter-feeding adaptations likely evolved independently from other extant filter-feeding sharks, even the lamniform basking shark, making it an example of convergent evolution.
Description
The appearance of the megamouth is distinctive, but little else is known about it. It has a brownish-black colour on top, is white underneath, and has an asymmetrical tail with a long upper lobe, similar to that of the thresher shark. The interior of its gill slits are lined with finger-like gill rakers that capture its food. A relatively poor swimmer, the megamouth has a soft, flabby body and lacks caudal keels. The megamouth is considerably less active than the other filter-feeding sharks, the basking shark and the whale shark. The megamouth has a stout body and a long, wide bulbous head.
Megamouths are large sharks, able to grow to in length. Mature males average at and females at . Weights of up to have been reported. A 2019 study suggested that it would have reached in maximum length. Megamouth sharks can be found as far northward as northern Japan; southern California (LACM 43745-1) and near Punta Eugenia, Baja California, and Hawaii. Megamouth sharks can be found at a depth of up to 1,000 m (3,280 ft). Megamouth sharks are dark blue, brownish-black, or gray above, lighter below; with a white band along the upper jaw; while the posterior margin of its fins are white.
As their name implies, megamouths have a large mouth with small teeth, and a broad, rounded snout, causing observers to occasionally mistake the megamouth for a young orca. The protruding inside of the upper lip is a brilliant silvery-white, which is very visible when the mouth is open. This lip was initially thought to be possibly embedded by luminous photophores when the first shark was examined in the early 1980s, which may act as a lure for plankton, while the team examining the second shark in the mid-1980s instead proposed that the lower lip might glow with the white band used as a reflector of sorts, but neither theory has been proven. In 2020, a study concluded that this species of shark does not in fact produce any light; the white band was found to merely be highly reflective of light. This white band is present in both sexes and could be either a feeding mechanism or possibly be used as a means of identifying other individuals of megamouth sharks. Their mouths can reach up to in width. Megamouth sharks have up to 50 rows of teeth in their upper-jaw and up to 75 rows of teeth in their lower-jaw.
Behavior
In 1990, a 4.9-m (16-foot) male megamouth shark was caught near the surface off Dana Point, California. This individual was eventually released with a small radio tag attached to its soft body. The tag relayed depth and time information over a two-day period. During the day, the shark swam at a depth around , but as the sun set, it would ascend and spend the night at depths between . Both day and night, its progress was very slow, around . This pattern of vertical migration is seen in many marine animals as they track the movement of plankton in the water column. The shark captured in March 2009 was reportedly netted at a depth of .
Reproduction
Reproduction is ovoviviparous, meaning that the young sharks develop in eggs that remain within the mother's body until they hatch. Tissue samples were obtained from twenty-seven megamouths caught in a two-year period off the Hualien coast (eastern Taiwan), and two caught in Baja California, Mexico, and samples taken from GenBank to perform a population genetic analyses of the megamouth shark; the results indicated no genetic diversity between populations found in different geographical locations, which indicates the species forms a single, highly migratory, interbreeding population.
| Biology and health sciences | Sharks | Animals |
318037 | https://en.wikipedia.org/wiki/Lamniformes | Lamniformes | The Lamniformes (, from Greek lamna "fish of prey") are an order of sharks commonly known as mackerel sharks (which may also refer specifically to the family Lamnidae). It includes some of the most familiar species of sharks, such as the great white as well as less familiar ones, such as the goblin shark and megamouth shark.
Members of the order are distinguished by possessing two dorsal fins, an anal fin, five gill slits, eyes without nictitating membranes, and a mouth extending behind the eyes. Species in two families of Lamniformes – Lamnidae and Alopiidae – are distinguished for maintaining a higher body temperature than the surrounding waters.
Members of the group include macropredators, generally of medium-large size, including the largest macropredatory shark ever, the extinct Otodus megalodon, as well as large planktivores.
Although some authors have argued that the Late Jurassic Palaeocarcharias should be considered the oldest known lamniform, this is disputed. The earliest unambiguous records of lamniformes are from the Early Cretaceous. Lamniformes underwent a major adaptive radiation during the Cretaceous and became prominent elements of oceanic ecosystems. They reached their highest diversity during the Late Cretaceous, but severely declined during the K-Pg extinction, before rebounding to a high but lower diversity peak during the Paleogene. Lamniformes have severely declined over the last 20 million years, with only 15 species alive today, compared to over 290 extant species in the Carcharhiniformes, which have evolved into medium and large body sizes during the same timeframe. The causes of the decline are uncertain, but are likely to have involved both biotic factors like competition and non-biotic factors like temperature and sea level.
Species
The order Lamniformes includes 10 families with 22 species, with a total of seven living families and 15 living species:
Order Lamniformes
Family Alopiidae Bonaparte, 1838 (thresher sharks)
Genus Alopias Rafinesque, 1810
Alopias pelagicus Nakamura, 1935 (pelagic thresher)
Alopias superciliosus R. T. Lowe, 1841 (bigeye thresher)
Alopias vulpinus (Bonnaterre, 1788) (common thresher)
Family †Anacoracidae Capetta, 1987 (extinct, Cretaceous period)
Genus †Squalicorax (crow sharks)
Genus †Scindocorax
Genus †Nanocorax
Genus †Ptychocorax
Family †Aquilolamnidae Vullo et al., 2021? (eagle sharks) (extinct, Late Cretaceous period)
Genus †Aquilolamna Vullo et al., 2021
†Aquilolamna milarcae Vullo et al., 2021
Family Cetorhinidae Gill, 1862
Genus Cetorhinus Blainville, 1816
Cetorhinus maximus (Gunnerus, 1765) (basking shark)
†Cetorhinus huddlestoni (Welton, 2014)
†Cetorhinus piersoni (Welton, 2015)
Genus †Keasius (Welton, 2013)
Family †Eoptolamnidae (extinct, Late Cretaceous period)
Genus †Eoptolamna
†Eoptolamna eccentrolopha
Genus †Leptostyrax
†Leptostyrax macrorhiza
Genus †Protolamna
†Protolamna sokolovi
†Protolamna borodini
†Protolamna carteri
†Protolamna compressidens
†Protolamna gigantea
†Protolamna roanokeensis
Family Lamnidae J. P. Müller and Henle, 1838 (mackerel sharks or white sharks)
Genus Carcharodon A. Smith, 1838
Carcharodon carcharias (Linnaeus, 1758) (great white shark)
†Carcharodon hubbelli Ehret, Macfadden, Jones, Devries, Foster & Salas-Gismondi, 2012 (Hubbell's white shark)
†Carcharodon caifassii Lawley, 1876
†Carcharodon carcharias-f Linnaeus, 1758
Genus Isurus Rafinesque, 1810
Isurus oxyrinchus Rafinesque, 1810 (shortfin mako)
Isurus paucus Guitart-Manday, 1966 (longfin mako)
Genus Lamna Cuvier, 1816
Lamna ditropis Hubbs & Follett, 1947 (salmon shark)
Lamna nasus (Bonnaterre, 1788) (porbeagle)
Family †Otodontidae Gluckman, 1964 (extinct, Late Cretaceous to Pliocene) (megatoothed sharks)
Genus †Cretalamna Gluckman, 1958
Genus †Otodus (Agassiz, 1843)
†Otodus obliquus (Agassiz, 1838)
†Otodus angustidens (Agassiz, 1843)
†Otodus chubutensis (Agassiz, 1843)
†Otodus megalodon (Agassiz, 1843) (megalodon)
†Otodus auriculatus (Jordan, 1923)
†Otodus sokolovi (Zhelezko and Kozlov, 1999)
†Otodus poseidoni (Zhelezko and Kozlov, 1999)
†Otodus minor (Giebel, 1943)
†Otodus hastalis (Lawley, 876)
†Otodus limhamnensis (Davis, 1890)
†Otodus debrayi (Leriche, 1906)
†Otodus naidini (Zhelezko in Zhelezko & Kozlov)
Genus †Megaselachus
†Megaselachus subauriculatus? (Glickman, 1964)
Genus †Megalolamna Shimada et al., 2016
Genus †Palaeocarcharodon Casieer, 1960
Genus †Kenolamna Siversson, 2017
Family Megachasmidae Taylor, Compagno & Struhsaker, 1983
Genus Megachasma Taylor, Compagno & Struhsaker, 1983
Megachasma pelagios Taylor, Compagno & Struhsaker, 1983 (megamouth shark)
Family Mitsukurinidae D. S. Jordan, 1898
Genus Mitsukurina D. S. Jordan, 1898
Mitsukurina owstoni D. S. Jordan, 1898 (goblin shark)
Family Odontaspididae Müller & Henle, 1839
Genus Carcharias Rafinesque, 1810
Carcharias taurus Rafinesque, 1810 (sand tiger shark)
Genus Odontaspis Agassiz, 1838
Odontaspis ferox (Risso, 1810) (smalltooth sand tiger)
Odontaspis noronhai (Maul, 1955) (bigeye sand tiger)
Family Pseudocarchariidae Compagno, 1973
Genus Pseudocarcharias Cadenat, 1963
Pseudocarcharias kamoharai (Matsubara, 1936) (crocodile shark)
Family †Cardabiodontidae (extinct, Late Cretaceous period)
Genus †Cardabiodon Siverson, 1999
†Cardabiodon ricki Siverson, 1999
†Cardabiodon venator Siverson and Lindgren, 2005
Genus †Dwardius Siverson, 1999
Genus †Parotodus? Cappetta, 1980
Family †Cretoxyrhinidae (extinct, Late Cretaceous period)
Genus †Cretoxyrhina Agassiz, 1843
†Cretoxyrhina vraconensis Zhelezko, 2000
†Cretoxyrhina denticulata Glückman, 1957
†Cretoxyrhina agassizensis Underwood and Cumbaa, 2010
†Cretoxyrhina mantelli Agassiz, 1843 (ginsu shark)
Family †Serratolamnidae
Genus †Serratolamna
Family †Ptychodontidae (extinct, Cretaceous period)
Genus †Ptychodus (16+ species)
Phylogeny
Below is a cladogram showing relationships within Lamniformes. The topology of extant families is based on Vella & Vella (2020) and the placements of Cretoxyrhinidae and Otodontidae are based on Ferrón (2017), Cooper (2020), and Greenfield (2022).
Sustainable consumption
In 2010, Greenpeace International added the shortfin mako shark (Isurus oxyrinchus) to its seafood red list.
| Biology and health sciences | Sharks | Animals |
318051 | https://en.wikipedia.org/wiki/Law%20of%20mass%20action | Law of mass action | In chemistry, the law of mass action is the proposition that the rate of a chemical reaction is directly proportional to the product of the activities or concentrations of the reactants. It explains and predicts behaviors of solutions in dynamic equilibrium. Specifically, it implies that for a chemical reaction mixture that is in equilibrium, the ratio between the concentration of reactants and products is constant.
Two aspects are involved in the initial formulation of the law: 1) the equilibrium aspect, concerning the composition of a reaction mixture at equilibrium and 2) the kinetic aspect concerning the rate equations for elementary reactions. Both aspects stem from the research performed by Cato M. Guldberg and Peter Waage between 1864 and 1879 in which equilibrium constants were derived by using kinetic data and the rate equation which they had proposed. Guldberg and Waage also recognized that chemical equilibrium is a dynamic process in which rates of reaction for the forward and backward reactions must be equal at chemical equilibrium. In order to derive the expression of the equilibrium constant appealing to kinetics, the expression of the rate equation must be used. The expression of the rate equations was rediscovered independently by Jacobus Henricus van 't Hoff.
The law is a statement about equilibrium and gives an expression for the equilibrium constant, a quantity characterizing chemical equilibrium. In modern chemistry this is derived using equilibrium thermodynamics. It can also be derived with the concept of chemical potential.
History
Two chemists generally expressed the composition of a mixture in terms of numerical values relating the amount of the product to describe the equilibrium state.
Cato Maximilian Guldberg and Peter Waage, building on Claude Louis Berthollet's ideas about reversible chemical reactions, proposed the law of mass action in 1864. These papers, in Danish, went largely unnoticed, as did the later publication (in French) of 1867 which contained a modified law and the experimental data on which that law was based.
In 1877 van 't Hoff independently came to similar conclusions, but was unaware of the earlier work, which prompted Guldberg and Waage to give a fuller and further developed account of their work, in German, in 1879. Van 't Hoff then accepted their priority.
1864
The equilibrium state (composition)
In their first paper, Guldberg and Waage suggested that in a reaction such as
A + B <=> A' + B'
the "chemical affinity" or "reaction force" between A and B did not just depend on the chemical nature of the reactants, as had previously been supposed, but also depended on the amount of each reactant in a reaction mixture. Thus the law of mass action was first stated as follows:
When two reactants, A and B, react together at a given temperature in a "substitution reaction," the affinity, or chemical force between them, is proportional to the active masses, [A] and [B], each raised to a particular power
.
In this context a substitution reaction was one such as {alcohol} + acid <=> {ester} + water. Active mass was defined in the 1879 paper as "the amount of substance in the sphere of action". For species in solution active mass is equal to concentration. For solids, active mass is taken as a constant. , a and b were regarded as empirical constants, to be determined by experiment.
At equilibrium, the chemical force driving the forward reaction must be equal to the chemical force driving the reverse reaction. Writing the initial active masses of A,B, A' and B' as p, q, p' and q' and the dissociated active mass at equilibrium as , this equality is represented by
represents the amount of reagents A and B that has been converted into A' and B'. Calculations based on this equation are reported in the second paper.
Dynamic approach to the equilibrium state
The third paper of 1864 was concerned with the kinetics of the same equilibrium system. Writing the dissociated active mass at some point in time as x, the rate of reaction was given as
Likewise the reverse reaction of A' with B' proceeded at a rate given by
The overall rate of conversion is the difference between these rates, so at equilibrium (when the composition stops changing) the two rates of reaction must be equal. Hence
...
1867
The rate expressions given in Guldberg and Waage's 1864 paper could not be differentiated, so they were simplified as follows. The chemical force was assumed to be directly proportional to the product of the active masses of the reactants.
This is equivalent to setting the exponents a and b of the earlier theory to one. The proportionality constant was called an affinity constant, k. The equilibrium condition for an "ideal" reaction was thus given the simplified form
[A]eq, [B]eq etc. are the active masses at equilibrium. In terms of the initial amounts reagents p,q etc. this becomes
The ratio of the affinity coefficients, k'/k, can be recognized as an equilibrium constant. Turning to the kinetic aspect, it was suggested that the velocity of reaction, v, is proportional to the sum of chemical affinities (forces). In its simplest form this results in the expression
where is the proportionality constant. Actually, Guldberg and Waage used a more complicated expression which allowed for interaction between A and A', etc. By making certain simplifying approximations to those more complicated expressions, the rate equation could be integrated and hence the equilibrium quantity could be calculated. The extensive calculations in the 1867 paper gave support to the simplified concept, namely,
The rate of a reaction is proportional to the product of the active masses of the reagents involved.
This is an alternative statement of the law of mass action.
1879
In the 1879 paper the assumption that reaction rate was proportional to the product of concentrations was justified microscopically in terms of the frequency of independent collisions, as had been developed for gas kinetics by Boltzmann in 1872 (Boltzmann equation). It was also proposed that the original theory of the equilibrium condition could be generalised to apply to any arbitrary chemical equilibrium.
The exponents α, β etc. are explicitly identified for the first time as the stoichiometric coefficients for the reaction.
Modern statement of the law
The affinity constants, k+ and k−, of the 1879 paper can now be recognised as rate constants. The equilibrium constant, K, was derived by setting the rates of forward and backward reactions to be equal. This also meant that the chemical affinities for the forward and backward reactions are equal. The resultant expression
is correct even from the modern perspective, apart from the use of concentrations instead of activities (the concept of chemical activity was developed by Josiah Willard Gibbs, in the 1870s, but was not widely known in Europe until the 1890s). The derivation from the reaction rate expressions is no longer considered to be valid. Nevertheless, Guldberg and Waage were on the right track when they suggested that the driving force for both forward and backward reactions is equal when the mixture is at equilibrium. The term they used for this force was chemical affinity. Today the expression for the equilibrium constant is derived by setting the chemical potential of forward and backward reactions to be equal. The generalisation of the law of mass action, in terms of affinity, to equilibria of arbitrary stoichiometry was a bold and correct conjecture.
The hypothesis that reaction rate is proportional to reactant concentrations is, strictly speaking, only true for elementary reactions (reactions with a single mechanistic step), but the empirical rate expression
is also applicable to second order reactions that may not be concerted reactions. Guldberg and Waage were fortunate in that reactions such as ester formation and hydrolysis, on which they originally based their theory, do indeed follow this rate expression.
In general many reactions occur with the formation of reactive intermediates, and/or through parallel reaction pathways. However, all reactions can be represented as a series of elementary reactions and, if the mechanism is known in detail, the rate equation for each individual step is given by the expression so that the overall rate equation can be derived from the individual steps. When this is done the equilibrium constant is obtained correctly from the rate equations for forward and backward reaction rates.
In biochemistry, there has been significant interest in the appropriate mathematical model for chemical reactions occurring in the intracellular medium. This is in contrast to the initial work done on chemical kinetics, which was in simplified systems where reactants were in a relatively dilute, pH-buffered, aqueous solution. In more complex environments, where bound particles may be prevented from disassociation by their surroundings, or diffusion is slow or anomalous, the model of mass action does not always describe the behavior of the reaction kinetics accurately. Several attempts have been made to modify the mass action model, but consensus has yet to be reached. Popular modifications replace the rate constants with functions of time and concentration. As an alternative to these mathematical constructs, one school of thought is that the mass action model can be valid in intracellular environments under certain conditions, but with different rates than would be found in a dilute, simple environment .
The fact that Guldberg and Waage developed their concepts in steps from 1864 to 1867 and 1879 has resulted in much confusion in the literature as to which equation the law of mass action refers. It has been a source of some textbook errors. Thus, today the "law of mass action" sometimes refers to the (correct) equilibrium constant formula,
and at other times to the (usually incorrect) rate formula.
Applications to other fields
In semiconductor physics
The law of mass action also has implications in semiconductor physics. Regardless of doping, the product of electron and hole densities is a constant at equilibrium. This constant depends on the thermal energy of the system (i.e. the product of the Boltzmann constant, , and temperature, ), as well as the band gap (the energy separation between conduction and valence bands, ) and effective density of states in the valence and conduction bands. When the equilibrium electron and hole densities are equal, their density is called the intrinsic carrier density as this would be the value of and in a perfect crystal. Note that the final product is independent of the Fermi level :
Diffusion in condensed matter
Yakov Frenkel represented diffusion process in condensed matter as an ensemble of elementary jumps and quasichemical interactions of particles and defects. Henry Eyring applied his theory of absolute reaction rates to this quasichemical representation of diffusion. Mass action law for diffusion leads to various nonlinear versions of Fick's law.
In mathematical ecology
The Lotka–Volterra equations describe dynamics of the predator-prey systems. The rate of predation upon the prey is assumed to be proportional to the rate at which the predators and the prey meet; this rate is evaluated as xy, where x is the number of prey, y is the number of predator. This is a typical example of the law of mass action.
In mathematical epidemiology
The law of mass action forms the basis of the compartmental model of disease spread in mathematical epidemiology, in which a population of humans, animals or other individuals is divided into categories of susceptible, infected, and recovered (immune). The principle of mass action is at the heart of the transmission term of compartmental models in epidemiology, which provide a useful abstraction of disease dynamics. The law of mass action formulation of the SIR model corresponds to the following "quasichemical" system of elementary reactions:
The list of components is S (susceptible individuals), I (infected individuals), and R (removed individuals, or just recovered ones if we neglect lethality);
The list of elementary reactions is
S + I -> 2I
I -> R.
If the immunity is unstable then the transition from R to S should be added that closes the cycle (SIRS model):
R -> S.
A rich system of law of mass action models was developed in mathematical epidemiology by adding components and elementary reactions.
Individuals in human or animal populations unlike molecules in an ideal solution do not mix homogeneously. There are some disease examples in which this non-homogeneity is great enough such that the outputs of the classical SIR model and their simple generalizations like SIS or SEIR, are invalid. For these situations, more sophisticated compartmental models or distributed reaction-diffusion models may be useful.
| Physical sciences | Kinetics | Chemistry |
318069 | https://en.wikipedia.org/wiki/Goblin%20shark | Goblin shark | The goblin shark (Mitsukurina owstoni) is a rare species of deep-sea shark. Sometimes called a "living fossil", it is the only extant representative of the family Mitsukurinidae, a lineage some 125 million years old. This pink-skinned animal has a distinctive profile with an elongated, flat snout, and highly protrusible jaws containing prominent nail-like teeth. It is usually between long when mature, though it can grow considerably larger such as one captured in 2000 that is thought to have measured . Goblin sharks are benthopelagic creatures that inhabit upper continental slopes, submarine canyons, and seamounts throughout the world at depths greater than , with adults found deeper than juveniles. Some researchers believe that these sharks could also dive to depths of up to , for short periods of time.
Various anatomical features of the goblin shark, such as its flabby body and small fins, suggest that it is sluggish in nature. This species hunts for teleost fishes, cephalopods, and crustaceans both near the sea floor and in the middle of the water column. Its long snout is covered with ampullae of Lorenzini that enable it to sense minute electric fields produced by nearby prey, which it can snatch up by rapidly extending its jaws. Small numbers of goblin sharks are unintentionally caught by deepwater fisheries. The International Union for Conservation of Nature (IUCN) has assessed it as Least Concern, despite its rarity, citing its wide distribution and low incidence of capture.
Taxonomy
American ichthyologist David Starr Jordan described the goblin shark in an 1898 issue of Proceedings of the California Academy of Sciences, recognizing the peculiar fish not only as a new species, but also a new genus and family. He based his account on an immature male caught in Sagami Bay near Yokohama, Japan. The specimen had been acquired by shipmaster and naturalist Alan Owston, who had given it to Professor Kakichi Mitsukuri at the University of Tokyo, who in turn brought it to Jordan. Jordan named the shark Mitsukurina owstoni in honor of these two men.
The common name "goblin shark" is a calque of its traditional Japanese name , a being a Japanese mythical creature often depicted with a long nose and red face. Another name for this species is elfin shark.
Soon after Jordan's description was published, several scientists noted the similarity between Mitsukurina and the extinct Mesozoic shark Scapanorhynchus. For a time, the prevailing opinion was to treat Mitsukurina as a junior synonym of Scapanorhynchus. Eventually, more complete fossils revealed many anatomical differences between Scapanorhynchus and Mitsukurina, causing modern authors to again regard them as distinct genera. Several goblin shark specimens were described as separate species from 1904 to 1937, none of which are now considered valid. This taxonomic confusion began because the specimens' jaws were fixed at varying degrees of protrusion during preservation, giving the appearance of proportional differences among the heads.
Phylogeny and evolution
Phylogenetic studies based on morphology have classified the goblin shark as the most basal member of the order Lamniformes, known as mackerel sharks. Studies using genetic data have also confirmed a basal classification for this species. The family Mitsukurinidae, represented by Mitsukurina, Scapanorhynchus, and Anomotodon, dates back to the Aptian age of the Cretaceous period (c. 125–113 Ma). Mitsukurina itself first appears in the fossil record during the period Middle Eocene (c. 49–37 Ma); extinct species include M. lineata and M. maslinensis. Striatolamia macrota, which lived in warm shallow waters during the Paleogene period (c. 66–23 Ma), may also be a Mitsukurina species. As the last member of an ancient lineage, and one that retains several "primitive" traits, the goblin shark has been described as a "living fossil".
Description
The goblin shark has a distinctively long and flat snout, resembling a blade. The proportional length of the snout decreases with age. The eyes are small and lack protective nictitating membranes; behind the eyes are spiracles. The large mouth is parabolic in shape. The jaws are very protrusible and can be extended almost to the end of the snout, though normally they are held flush against the underside of the head. It has 35–53 upper and 31–62 lower tooth rows. The teeth in the main part of the jaws are long and narrow, particularly those near the symphysis (jaw midpoint), and are finely grooved lengthwise. The rear teeth near the corners of the jaw are small and have a flattened shape for crushing. Much individual variation of tooth length and width occurs, as for whether the teeth have a smaller cusplet on each side of the main cusp, and regarding the presence of toothless gaps at the symphysis or between the main and rear teeth. The five pairs of gill slits are short, with the gill filaments inside partly exposed; the fifth pair is above the origin of the pectoral fins.
The body is fairly slender and flabby. The two dorsal fins are similar in size and shape, both being small and rounded. The pectoral fins are also rather small and rounded. The pelvic and anal fins have long bases and are larger than the dorsal fins. The caudal peduncle is flattened from side-to-side and lacks keels or notches. The asymmetric caudal fin has a long upper lobe with a shallow ventral notch near the tip, and an indistinct lower lobe. The soft, semi-translucent skin has a rough texture from a covering of dermal denticles, each shaped like a short upright spine with lengthwise ridges. Living sharks of this species are pink or tan due to visible blood vessels beneath the skin; the color deepens with age, and young sharks may be almost white. The fins' margins are translucent gray or blue, and the eyes are black with bluish streaks in the irises. After death, the coloration fades quickly to dull gray or brown. Adult sharks usually measure between long. However, the capture of an enormous female estimated at long during 2000 showed this species can grow far larger than suspected previously. A 2019 study suggested that it would have reached in maximum length. Until 2022, the maximum weight recorded was for a shark of 3.8 m (12.5 ft) in length. In 2023, a heavily pregnant, individual weighing , was landed in Taiwan. The enormous individual sparked criticism of the fishing method of bottom trawling which was used to catch it.
Distribution and habitat
The goblin shark has been caught in all three major oceans, indicating a wide global distribution. In the Atlantic Ocean, it has been recorded from the northern Gulf of Mexico, Suriname, French Guiana, and southern Brazil in the west, and France, Portugal, Madeira, and Senegal in the east. It has also been collected from seamounts along the Mid-Atlantic Ridge. In the Indo-Pacific and Oceania, it has been found off South Africa, Mozambique, Japan, Taiwan, Australia and New Zealand. This species has been recorded from off East Cape to Kaikōura Canyon and from the Challenger Plateau near New Zealand. A single eastern Pacific specimen is known, collected off southern California. This species is most often found over the upper continental slope at depths of . It has been caught as deep as , a tooth has been found lodged in an undersea cable at a depth of , and a live individual was filmed at a depth of in the Tonga Trench. Adults inhabit greater depths than juveniles. Immature goblin sharks frequent the submarine canyons off southern Japan at depths of , with individuals occasionally wandering into inshore waters as shallow as . Two juvenile goblin sharks were captured from the Tokyo Underwater Canyon in Tokyo Bay, filmed in the Kanaya Fishing Port, and released on 30 January 2008 and 24 January 2011. The goblin shark filmed in 2008 was caught at a depth of 150–350m (492–1,148 ft). On 19 April 2014, fishermen in Key West, Florida, while fishing in the Gulf of Mexico, caught a goblin shark in their fishing net, only the second one ever to be caught in the Gulf. The shark was photographed and released back into the water. The first shark found in the Gulf was caught by commercial fisherman on 25 July 2000 at a depth of approximately 919–1,099 m (3,016–3,606 ft) and is thought to have been about 20 ft long. During July 2014, a goblin shark was found in a fishery net in Sri Lanka, near the eastern coast of Sri Lanka. The shark was about long and weighed about . The shark was given to the NARA (National Aquatic Resource Research & Development Agency) for further research.
Biology and ecology
Although observations of living goblin sharks are scant, its anatomy suggests its lifestyle is inactive and sluggish. Its skeleton is reduced and poorly calcified, the muscle blocks along its sides (myomeres) are weakly developed, and its fins are soft and small. Its long caudal fin, held at a low angle, is also typical of a slow-swimming shark. The long snout appears to have a sensory function, as it bears numerous ampullae of Lorenzini that can detect the weak electric fields produced by other animals. Due to the snout's softness, it is unlikely to be used for stirring up prey from the bottom as has been proposed. Vision seems to be less important than other senses, considering the relatively small optic tectum in the shark's brain. Yet unlike most deep-sea sharks, it can change the size of its pupils, thus probably does use its sight in some situations. Goblin sharks may be the prey of blue sharks (Prionace glauca). Parasites documented from this species include the copepod Echthrogaleus mitsukurinae, and the tapeworms Litobothrium amsichensis and Marsupiobothrium gobelinus.
Feeding
The goblin shark feeds mainly on teleost fishes such as rattails and dragonfishes. It also consumes cephalopods and crustaceans, including decapods and isopods. Garbage has been recorded from the stomachs of some specimens. Its known prey includes bottom-dwelling species such as the blackbelly rosefish (Helicolenus dactylopterus), and midwater species such as the squid Teuthowenia pellucida and the ostracod Macrocypridina castanea rotunda. Thus, the goblin shark appears to forage for food both near the sea floor and far above it.
Since it is not a fast swimmer, the goblin shark may be an ambush predator. Its low-density flesh and large oily liver make it neutrally buoyant, allowing it to drift towards its prey with minimal motions so as to avoid detection. Once prey comes into range, the shark's specialized jaws can snap forward to capture it. The protrusion of the jaw is assisted by two pairs of elastic ligaments associated with the mandibular joint, which are pulled taut when the jaws are in their normal retracted position; when the shark bites, the ligaments release their tension and essentially "catapult" the jaws forward. At the same time, the well-developed basihyal (analogous to a tongue) on the floor of the mouth drops, expanding the oral cavity and sucking in water and prey. Striking and prey capture events were videotaped and recorded for the first time during 2008 and 2011 and helped to confirm the use and systematics of the protrusible jaws of goblin sharks. The video evidence suggests that while the jaws are definitely unique, goblin sharks use ram feeding, a type of prey capture that is typical of many mackerel sharks. What makes the goblin shark unique is the kinematics of their jaw when feeding. The lower jaw seems to undergo more complex movements and is important in capturing the prey. The measured protrusions of the upper and lower jaw combined put the goblin shark jaws at 2.1–9.5 times more protrusible than other sharks. The lower jaw has a velocity about two times greater than the upper jaw because it not only protrudes forward, but also swings upward to capture the prey, and the maximum velocity of the jaws is 3.14 m/s. The goblin shark has a re-opening and re-closing pattern during the strike, a behavior that has never been seen in other sharks before and could be related to the extent with which the goblin shark protrudes its jaws.
Growth and reproduction
The reproductive behaviors of the goblin shark are poorly understood. Mating has never been observed between two goblin sharks, and a pregnant female has yet to be captured. It likely shares the reproductive characteristics of other mackerel sharks, which are viviparous with small litter sizes and embryos that grow during gestation by eating undeveloped eggs (oophagy). The birth size is probably close to , the length of the smallest known specimen. Males mature sexually at about long, while female maturation size is unknown. No data is available concerning growth and aging. Some researchers have estimated, based on their own research and prior findings, that male goblin sharks mature at approximately 16 years old and can live up to 60 years.
Human interactions
Given the depths at which it lives, the goblin shark poses little danger to humans. The first known findings pertaining to the goblin shark were published in 1910. "The new shark is certainly grotesque, [...] the most remarkable feature is the curiously elongated nose." A few specimens have been collected alive and brought to public aquariums, though they only survived briefly. One was kept at Tokai University and lived for a week, while another was kept at Tokyo Sea Life Park and lived for two days. Its economic significance is minimal; the meat may be dried and salted, while the jaws fetch high prices from collectors. At one time, the Japanese also used it for liver oil and fertilizer. This shark is not targeted by any fisheries, but is occasionally found as bycatch in bottom gillnets and trawls, hooked on longlines, or entangled in fishing gear. Most captures are isolated incidents; one of the few areas where it is caught regularly is off southern Japan, where around 30 individuals (mostly juveniles) are taken each year. A black scabbardfish (Aphanopus carbo) fishery off Madeira also takes two or three goblin sharks annually. During April 2003, more than a hundred goblin sharks were caught off northwestern Taiwan; the cause of the event was unknown, though observers noted it was preceded by a major earthquake. The species had never been recorded in the area before, nor has it been found in such numbers since. The International Union for Conservation of Nature (IUCN) has categorized the goblin shark as Least Concern. In addition to its wide range, most of its population is thought to reside in unfished environments because few adults are caught. Therefore, it is not believed to be threatened by human activity. However, during June 2018 the New Zealand Department of Conservation classified the goblin shark as "At Risk – Naturally Uncommon" with the qualifiers "Data Poor" and "Secure Overseas" using the New Zealand Threat Classification System.
| Biology and health sciences | Sharks | Animals |
318200 | https://en.wikipedia.org/wiki/Giraffidae | Giraffidae | The Giraffidae are a family of ruminant artiodactyl mammals that share a recent common ancestor with deer and bovids. This family, once a diverse group spread throughout Eurasia and Africa, presently comprises only two extant genera, the giraffe (between one and eight, usually four, species of Giraffa, depending on taxonomic interpretation) and the okapi (the only known species of Okapia). Both are confined to sub-Saharan Africa: the giraffe to the open savannas, and the okapi to the dense rainforest of the Congo. The two genera look very different on first sight, but share a number of common features, including a long, dark-coloured tongue, lobed canine teeth, and horns covered in skin, called ossicones.
Taxonomy
Evolutionary background
The giraffids are ruminants of the clade Pecora. Other extant pecorans are the families Antilocapridae (pronghorns), Cervidae (deer), Moschidae (musk deer), and Bovidae (cattle, goats and sheep, wildebeests and allies, and antelopes). The exact interrelationships among the pecorans have been debated, mainly focusing on the placement of Giraffidae, but a recent large-scale ruminant genome sequencing study suggests Antilocapridae are the sister taxon to Giraffidae, as shown in the cladogram below.
The ancestors of pronghorn diverged from the giraffids in the Early Miocene. This was in part of a relatively late mammal diversification following a climate change that transformed subtropical woodlands into open savannah grasslands.
The fossil record of giraffids and their stem-relatives is quite intensive, with fossil of these taxa include Gelocidae, Palaeomerycidae, Prolibytheridae, and Climacoceratidae. It is thought that the palaeomerycids, prolibytherids, climacoceratids and the giraffids all form a clade of pecorans known as Giraffomorpha. The relationship between the climacoceratids and giraffids is supported by the presence of a bilobed canine, and have been postulated into two hypotheses. One is the climacoceratids were the ancestors of the sivatheres, as both groups were large, deer-like giraffoids with branching antler-like ossicones, while an extinct basal group of giraffoids, canthumerycines, evolved into the ancestors of Giraffidae. Another more commonly supported hypothesis is climacoceratids were merely the sister clade to giraffids, with sivatheres being either basal giraffids or descended from a lineage that also includes the okapi. While the current range of giraffids today is in Africa, the fossil record of the group has shown this family was once widespread throughout of Eurasia.
Below is the phylogenetic relationships of giraffomorphs after Solounias (2007), Sánchez et al. (2015) and Ríos et al. (2017):
Classification
Below is the total taxonomy of valid extant and fossil taxa (as well as junior synonyms which are listed in the brackets).
Family Giraffidae J.E.Gray, 1821
Basal extinct giraffids
†Csakvarotherium Kretzoi, 1930
†Csakvarotherium hungaricum Kretzoi, 1930
†Injanatherium Heintz, Brunet & Sen, 1981
†Injanatherium arabicum Morales, Soria & Thomas, 1987
†Injanatherium hazimi Heintz, Brunet & Sen, 1981
†Propalaeomeryx Lydekker, 1883 [Progiraffa Pilgrim, 1908]
†Propalaeomeryx sivalensis Lydekker, 1883 [Progiraffa exigua Pilgrim, 1908]
†Shansitherium Killgus, 1922 [Schansitherium [sic]]
†Shansitherium quadricornis (Bohlin, 1926) [Palaeotragus quadricornis Bohlin, 1926]
†Shansitherium tafeli Killgus, 1922
†Umbrotherium Abbazzi, Delfino, Gallai, Trebini & Rook, 2008
†Umbrotherium azzarolii Abbazzi, Delfino, Gallai, Trebini & Rook, 2008
Subfamily †Canthumerycinae Hamilton, 1978
†Georgiomeryx Paraskevaidis, 1940
†Georgiomeryx georgalasi Paraskevaidis, 1940
†Canthumeryx Hamilton 1973 [Zarafa Hamilton, 1973]
†Canthumeryx sirtensis Hamilton 1973 [Zarafa zelteni Hamilton, 1973]
Subfamily Giraffinae J.E.Gray, 1821
Tribe Giraffini J.E.Gray, 1821
Subtribe Giraffina J.E.Gray, 1821
Giraffa Brisson, 1762 [Camelopardalis von Schreber, 1784 and Orasius Oken, 1816]
Giraffa camelopardalis super-complex (Linnaeus, 1758)
Giraffa giraffa complex (von Schreber, 1784)
Giraffa angolensis Lydekker, 1903 – Angolan giraffe
Giraffa giraffa (von Schreber, 1784) – South African giraffe
Giraffa tippelskirchii complex Matschie, 1898
Giraffa thornicrofti Lydekker, 1911 – Rhodesian giraffe
Giraffa tippelskirchii Matschie, 1898 – Masai giraffe
Giraffa reticulata de Winton, 1899 – Reticulated giraffe
Giraffa camelopardalis complex (Linnaeus, 1758)
Giraffa peralta Thomas, 1898 – West African giraffe
Giraffa antiquorum Jardine & Swainson, 1835 – Kordofan giraffe
Giraffa camelopardalis (Linnaeus, 1758) – Northern giraffe
Giraffa camelopardalis rothschildi Lydekker, 1903 – Rothschild's giraffe
Giraffa camelopardalis camelopardalis (Linnaeus, 1758) – Nubian giraffe
†Giraffa jumae Leakey, 1967
†Giraffa priscilla Pilgrim, 1911
†Giraffa punjabiensis Pilgrim, 1911
†Giraffa pygmaea Harris, 1976
†Giraffa sivalensis (Falconer & Cautley, 1843) [Camelopardalis sivalensis Falconer & Cautley, 1843 and Camelopardalis affinis Falconer & Cautley, 1843]
†Giraffa stillei (Dietrich, 1942) [Okapia stillei Dietrich, 1942 and Giraffa gracilis Arambourg, 1947]
Subtribe †Bohlinina Solounias, 2007
†Bohlinia Matthew, 1929
†Bohlinia adoumi Likius, Vignaud & Brunet, 2007
†Bohlinia attica (Gaudry & Lartet, 1856) [Giraffa attica (Gaudry & Lartet, 1856) and Orasius attica (Gaudry & Lartet, 1856)]
†Bohlinia nikitiae Kostopoulos, Koliadimou & Koufos, 1996
†Honanotherium Bohlin, 1927
†Honanotherium bernori Solounias & Danowitz, 2016
†Honanotherium schlosseri (Pilgrim, 1911) [Giraffa schlosseri Pilgrim, 1911]
Tribe Palaeotragini Pilgrim, 1910
Subtribe †Palaeotragina Pilgrim, 1910
†Giraffokeryx Pilgrim, 1910
†Giraffokeryx anatoliensis Geraads & Aslan, 2003
†Giraffokeryx primaevus (Churcher, 1970) [Palaeotragus primaevus Churcher, 1970; Samotherium africanum Churcher, 1970 and Amotherium africanum [sic]]
†Giraffokeryx punjabiensis Pilgrim, 1910
†Mitilanotherium
†Mitilanotherium inexpectatum
†Palaeogiraffa Bonis & Bouvrain, 2003
†Palaeogiraffa macedoniae (Geraads, 1989) [Decennatherium macedoniae Geraads, 1989]
†Palaeogiraffa major Bonis & Bouvrain, 2003
†Palaeogiraffa pamiri (Ozansoy, 1965) [Samotherium pamiri Ozansoy, 1965]
†Palaeotragus Gaudry, 1861 [Achtiaria Borissiak, 1914; Macedonitherium Sickenberg, 1967; Mitilanotherium Samson & Radulesco, 1966 and Sogdianotherium Sharapov, 1974]
†Palaeotragus coelophrys (Rodler & Weithofer, 1890) [Alcicephalus coelophrys Rodler & Weithofer, 1890]
†Palaeotragus germaini Arambourg, 1959
†Palaeotragus inexspectatus (Samson & Radulesco, 1966) [Macedonitherium martinii Sickenberg, 1967; Mitilanotherium inexpectatum Samson & Radulesco, 1966; Mitilanotherium kuruksaense (Sharapov, 1974); Mitilanotherium martinii (Sickenberg, 1967); Palaeotragus priasovicus Godina & Bajgusheva, 1985 and Sogdianotherium kuruksaense Sharapov, 1974]
†Palaeotragus lavocanti Heintz, 1976
†Palaeotragus robinsoni Crusafont-Pairó, 1979
†Palaeotragus rouenii Gaudry, 1861 [Palaeotragus microdon Koken, 1885]
†Palaeotragus tungurensis Colbert, 1936
†Praepalaeotragus Godina, Vislobokova & Abdrachmanova, 1993
†Praepalaeotragus actaensis Godina, Vislobokova & Abdrachmanova, 1993
†Samotherium Forsyth Major, 1888 [Alcicephalus Rodler & Weithofer, 1890; Chersenotherium Alexajew, 1916 and Amotherium [sic]]
†Samotherium boissieri Forsyth Major, 1888
†Samotherium eminens (Alexajew, 1916) [Chersenotherium eminens Alexajew, 1916]
†Samotherium major Bohlin, 1926
†Samotherium neumayri (Rodler & Weithofer, 1890) [Alcicephalus neumayri Rodler & Weithofer, 1890]
†Samotherium sinense (Schlosser, 1903) [Alcicephalus sinense Schlosser, 1903]
Subtribe Okapiina Bohlin, 1926
†Afrikanokeryx Harris, Solounias & Geraads, 2010
†Afrikanokeryx leakeyi Harris, Solounias & Geraads, 2010
Okapia Lankester, 1901
Okapia johnstoni (P. L. Sclater, 1901) – Okapi
†Subfamily Sivatheriinae Bonaparte, 1850
†Birgerbohlinia Crusafont Pairó, 1952
†Birgerbohlinia schaubi Crusafont Pairó, 1952
†Bramatherium Falconer, 1845 [Hydaspitherium Lydekker, 1876]
†Bramatherium giganteus Khan & Sarwar, 2002
†Bramatherium grande (Lydekker, 1878) [Hydaspitherium grande Lydekker, 1878]
†Bramatherium magnum (Pilgrim, 1910) [Hydaspitherium magnum Pilgrim, 1910]
†Bramatherium megacephalum (Lydekker, 1876) [Hydaspitherium megacephalum Lydekker, 1876]
†Bramatherium perimense Falconer, 1845
†Bramatherium progressus Khan, Sarwar & Khan, 1993
†Bramatherium suchovi Godina, 1977
†Decennatherium Crusafont Pairó, 1952
†Decennatherium rex Ríos, Sánchez & Morales, 2017
†Decennatherium pachecoi Crusafont Pairó, 1952
†Helladotherium Gaudry, 1860
†Helladotherium duvernoyi (Gaudry & Lartet, 1856) [Camelopardalis duvernoyi Gaudry & Lartet, 1856]
†Sivatherium Falconer & Cautley, 1836 [Griquatherium Haughton, 1922; Indratherium Pilgrim, 1910; Libytherium Pomel, 1892 and Orangiatherium van Hoepen, 1932]
†Sivatherium giganteum Falconer & Cautley, 1836
†Sivatherium hendeyi Harris, 1976
†Sivatherium maurusium (Pomel, 1892) [Libytherium maurusium Pomel, 1892; Griquatherium cingulatum Haughton, 1922; Helladotherium olduvaiense Hopwood, 1934; Sivatherium olduvaiense (Hopwood, 1934); Libytherium olduvaiense Hopwood, 1934 and Orangiatherium vanrhyni van Hoepen, 1932]
†Vishnutherium Lydekker, 1876
†Vishnutherium iravadicum Lydekker 1876
Characteristics
The giraffe stands tall, with males taller than females. The giraffe and the okapi have characteristic long necks and long legs. Ossicones are present on males and females in the giraffe, but only on males in the okapi.
Giraffids share many common features with other ruminants. They have cloven hooves and cannon bones, much like bovids, and a complex, four-chambered stomach. They have no upper incisors or upper canines, replacing them with a tough, horny pad. An especially long diastema is seen between the front and cheek teeth. The latter are selenodont, adapted for grinding up tough plant matter. Like most other ruminants, the dental formula for giraffids is . Giraffids have prehensile tongues (specially adapted for grasping).
The extant giraffids, the forest-dwelling okapi and the savannah-living giraffe, have several features in common, including a pair of skin-covered horns, called ossicones, up to long (absent in female okapis); a long, black, prehensile tongue; lobed canine teeth; patterned coats acting as camouflage; and a back sloping towards the rear. The okapi's neck is long compared to most ruminants, but not nearly so long as the giraffe's. Male giraffes are the tallest of all mammals: their horns reach above the ground and their shoulder , whereas the okapi has a shoulder height of .
Distribution
The two extant genera are now confined to sub-Saharan Africa. The okapi is restricted to a small range in the northern rainforest of the Democratic Republic of Congo. Although the range of the giraffe is considerably larger, it once covered an area twice the present size – all parts of Africa that could offer an arid and dry landscape furnished with trees.
Behavior
The social structure and behavior is markedly different in okapis and giraffes, but although little is known of the okapi's behavior in the wild, a few things are known to be present in both species:
They have an ambling gait similar to camels, with their weight supported alternately by their left and right legs, while their necks maintain balance. Giraffes can run up to this way and are documented to have covered in the Sahel during the dry season.
The dominance hierarchy, which has been well-documented among giraffes, has also been seen among captive okapis. An adult giraffe head can weigh , and if necessary, male giraffes establish a hierarchy among themselves by swinging their heads at each other, horns first, a behavior known as "necking". A subordinate okapi signals submission by placing its head and neck on the ground.
Giraffes are sociable, whereas okapis live mainly solitary lives. Giraffes temporarily form herds of up to 20 individuals; these herds can be mixed or uniform groups of males and females, young and adults. Okapis are normally seen in mother-offspring pairs, although they occasionally gather around a prime food source. Giraffe are not territorial, but have ranges that can dramatically vary between – – depending on food availability, whereas okapis have individual ranges about in size.
Giraffes and okapis are normally silent, but both have a range of vocalizations, including coughing, snorting, moaning, hissing, and whistling. Giraffes have been suggested to be able to communicate using infrasonic sounds like elephants and blue whales.
| Biology and health sciences | Giraffidae | Animals |
318215 | https://en.wikipedia.org/wiki/African%20penguin | African penguin | The African penguin (Spheniscus demersus), also known as Cape penguin or South African penguin, is a species of penguin confined to southern African waters. It is the only penguin found in the Old World. Like all penguins, it is flightless, with a streamlined body and wings stiffened and flattened into flippers for a marine habitat. Adults weigh an average of and are tall. The species has distinctive pink patches of skin above the eyes and a black facial mask. The body's upper parts are black and sharply delineated from the white underparts, which are spotted and marked with a black band.
The African penguin is a pursuit diver and feeds primarily on fish and squid. Once extremely numerous, the African penguin is now the rarest species of penguin classified as critically endangered, with its population declining rapidly due to a combination of several threats. It is a charismatic species and is popular with tourists. Other vernacular names of the species include black-footed penguin and jackass penguin, due to the species' loud, donkey-like noise (although several related species of South American penguins produce the same sound). They can be found along the coast of South Africa and Namibia.
Taxonomy
In 1747, the English naturalist George Edwards included an illustration and a description of the African penguin in the second volume of his A Natural History of Uncommon Birds. He used the English name "The Black-Footed Penguins". Edwards based his hand-coloured etching on two preserved specimens that had been brought to London. He suspected that they had been collected near the Cape of Good Hope. When in 1758 the Swedish naturalist Carl Linnaeus updated his Systema Naturae for the tenth edition, he placed the African penguin with the wandering albatross in the genus Diomedea. Linnaeus included a brief description, coined the binomial name Diomedea demersa and cited Edwards' work. The African penguin is now placed with the banded penguins in the genus Spheniscus that was introduced in 1760 by the French zoologist Mathurin Jacques Brisson. The genus name Spheniscus is from Ancient Greek word σφήν (sphēn) meaning "wedge" and is a reference to the animal's thin, wedge-shaped flippers. The specific epithet demersus is Latin meaning "plunging" (from demergere meaning "to sink").
Banded penguins are found mainly in the Southern Hemisphere with the Humboldt penguin and Magellanic penguins found in southern South America and the Galápagos penguin found in the Pacific Ocean near the equator. All are similar in shape, colour and behaviour.
Description
African penguins grow to tall and weigh between . The beak length of the African penguin varies, usually growing between . They have a black stripe and black spots on the chest, the pattern of which is unique to each penguin, like human fingerprints. The sweat glands above the eyes cool the birds' blood and as the temperature rises, increased blood flow causes the glands to get pinker. This species exhibits slight sexual dimorphism; the males are slightly larger than the females and have longer beaks. Juveniles do not possess the bold, delineated markings of the adult, but instead have dark upperparts that vary from greyish-blue to brown; the pale underparts lack both spots and the band. The beak is more pointed than that of the Humboldt penguin. The African penguin's colouring is a form of protective colouration known as countershading. The white undersides of the birds are difficult to spot by predators under the water and the penguins' black backs blend in with the water when viewed from above.
African penguins resemble and are thought to be related to the Humboldt, Magellanic and Galápagos penguins. African penguins have a very recognisable appearance, with a thick band of black that is in the shape of an upside-down horseshoe. They have black feet and black spots that vary in size and shape between individuals. Magellanic penguins share a similar bar marking that often confuses the two; the Magellanic has a double bar on the throat and chest, whereas the African has a single bar. These penguins have the nickname "jackass penguin", which comes from the loud penguin noises they make.
Distribution and habitat
The African penguin is only found on the southwestern coast of Africa, living in colonies on 24 islands between Namibia and Algoa Bay, near Port Elizabeth, South Africa. It is the only penguin species that breeds in Africa and its presence gave name to the Penguin Islands.
Two colonies were established by penguins in the 1980s on the mainland near Cape Town, namely Boulders Beach near Simon's Town and Stony Point in Betty's Bay. Mainland colonies likely became possible only in recent times due to the reduction of predator numbers, although the Betty's Bay colony has been attacked by leopards. The only other mainland colony is in Namibia, but it is not known when it was established.
Boulders Beach is a tourist attraction due to the beach, swimming and the penguins. The penguins will allow people to approach them as close as a meter.
Breeding populations of African penguins are being kept in numerous zoos worldwide. No colonies are known outside the southwestern coast of Africa, although vagrants (mostly juveniles) may occasionally be sighted beyond the normal range.
Population
Roughly 4 million African penguins existed at the beginning of the 19th century. Of the 1.5 million African penguins estimated in 1910, only some 10% remained at the end of the 20th century. African penguin populations, which breed in Namibia and South Africa, have declined by 95% since pre-industrial times.
Today, their breeding is largely restricted to 24 islands from Namibia to Algoa Bay, South Africa, with the Boulders Beach colony being an exception to this rule.
The total population fell to approximately 150,000–180,000 in 2000. Of those, 56,000 belonged to the Dassen Island colony and 14,000 to the Robben Island colony. The colony at Dyer Island in South Africa fell from 46,000 in the early 1970s to 3,000 in 2008.
In 2008, 5,000 breeding pairs were estimated to live in Namibia.
In 2010, the total African penguin population was estimated at 55,000. At the rate of decline seen from 2000 to 2010, the African penguin is expected to be extinct in the wild by 2026.
In 2012, about 18,700 breeding pairs were estimated to live in South Africa, with the majority on St. Croix Island in Algoa Bay.
The total breeding population across both South Africa and Namibia fell to a historic low of about 20,850 pairs in 2019.
Behaviour
Diet
African penguins forage in the open sea, where they feed on pelagic fish such as sardines (including the blue pilchard), Cape horse mackerels, red-eye round herrings and anchovies (specifically the European anchovy and the Southern African anchovy) and marine invertebrates such as squids and small crustaceans, primarily krills and shrimps. Penguins normally swim within of the shore. A penguin may consume up to of prey every day, but this may increase to over when raising older chicks.
Due to the marked decline of sardines in the waters near its habitat, African penguins' diet has shifted towards anchovies to some extent, although available sardine biomass is still a notable determinant of penguin population development and breeding success. While a diet of anchovies appears to be generally sufficient for the penguins, it is not ideal due to anchovies' lower concentrations of fat and protein. The Penguin's diet changes throughout the year; as in many seabirds, it is believed that the interaction of diet choice and breeding success helps the penguins maintain their population size. Although parent penguins are protective of their chicks, they will not incur nutritional deficits themselves if food is scarce and hunting requires a greater time or energy commitment. This may lead to higher rates of brood loss under poor food conditions.
When foraging, African penguins carry out dives that reach an average depth of and last for 69 seconds, although a maximum depth of and duration of 275 seconds has been recorded.
Breeding
The African penguin is monogamous; it breeds in colonies and pairs return to the same site each year. It has an extended breeding season, with nesting usually peaking from March to May in South Africa and November to December in Namibia. A clutch of two eggs is laid either in burrows burrowed in guano or nests in the sand under boulders or bushes. Incubation is undertaken equally by both parents for around 40 days. At least one parent guards the chicks for about one month, whereafter the chicks join a crèche with other chicks and both parents spend most of the day foraging in the sea.
Chicks fledge at 60 to 130 days, the timing depending on environmental factors such as the quality and availability of food. The fledged chicks then go to sea on their own, where they spend the next one to nearly two years. They then return to their natal colony to moult into adult plumage.
When penguins moult, they are unable to forage in the sea as their new feathers are not yet waterproof; therefore, they fast over the entire moulting period. African penguins typically take around three weeks to moult and lose about half of their body weight by burning up their fat reserves in the process.
African penguins spend most of their lives at sea until it comes time for them to lay their eggs. Females remain fertile for about 10 years. Due to high predation on the mainland, African penguins will seek protection on offshore islands, where they are safer from larger mammals and natural challenges. These penguins usually breed during the winter when temperatures are cooler. African penguins often will abandon their eggs if they become overheated in the hot sun and abandoned eggs never survive the heat. The eggs are three to four times bigger than chicken eggs. Ideally, the eggs are incubated in a burrow dug into the guano layer (which provides suitable temperature regulation), but the widespread human removal of guano deposits has rendered this type of nest unfeasible in many colonies. To compensate, penguins burrow holes in the sand, nest under rocks or bushes or make use of nest boxes if they are provided. The penguins spend three weeks on land caring for their offspring, after which chicks may be left alone during the day while the parents forage. The chicks are frequently killed by predators or succumb to the hot sun. Parents usually feed hatchlings during dusk or dawn.
In 2015, when foraging conditions were favourable, more male than female African penguin chicks were produced in the colony on Bird Island. Male chicks also had higher growth rates and fledging mass and therefore may have higher post-fledging survival than females. This, coupled with higher adult female mortality in this species, may result in a male-biased adult sex ratio and may indicate that conservation strategies focused on benefiting female African penguins may be necessary.
Predation
The average lifespan of an African penguin is 10 to around 25 years in the wild and up to 30 in captivity.
The primary predators of African penguins at sea include sharks and fur seals. While nesting: kelp gulls, Cape genets, mongooses, caracals and domestic cats and dogs may prey on the penguins and their chicks. Mortality from terrestrial predators is higher if penguins are forced to breed in the open in the absence of suitable burrows or nest boxes.
Threats and conservation
Historical exploitation
African penguin eggs were considered a delicacy and were still being eaten and collected for sale as recently as the 1970s. In the 1950s, they were being collected from Dassen Island and sold in nearby towns. In 1953, 12,000 eggs were collected. In the late 1950s, some French chefs expressed interest in recipes including African penguin eggs collected from the islands off the west coast of South Africa and placed annual orders for small quantities. In the mid-1960s, eggs were collected in the thousands and sold by the dozen, with each customer limited to two dozen eggs in total.
The practice of collecting African penguin eggs involved smashing those found a few days before a collecting effort to ensure that only freshly laid eggs were sold. This added to the drastic decline of the African penguin population around the Cape coast, a decline which was hastened by the removal of guano from islands for use as fertiliser, eliminating the burrowing material used by penguins.
Oil spills
Penguins remain susceptible to pollution of their habitat by petrochemicals from spills, shipwrecks and cleaning of tankers while at sea. Accounts of African penguins impacted by oil date back to the 1930s. African penguins' exposure to oil spills is both chronic (higher frequency small discharges of oil at sea) and acute (rare maritime disasters where large volumes of oil are released in a single event). Penguins of many species have been impacted by oil spills across the southern hemisphere.
In 1948, the tanker Esso Wheeling sank, subsequently oiling and killing thousands of penguins of the Dyer Island colony. In 1953, dead penguins were among a range of dead birds, fish and other marine life that washed ashore after the tanker Sliedrecht was holed and spilled oil near Table Bay. In 1971, the SS Wafra oil spill impacted the African penguin colony of Dyer Island. In 1972, oil spilt following the Oswego-Guardian and Texanita collision oiled roughly 500 penguins. In 1975, newspapers reported that oil pollution from shipwrecks and the pumping of bilges at sea had killed tens of thousands of African penguins. At the time, the Dassen Island colony was being passed by 650 oil tankers each month because the Suez Canal had become blocked with wrecked vessels, thus increasing maritime traffic past the Cape of Good Hope.
In 1979, an oil spill prompted the collection and treatment of 150 African penguins from St. Croix Island near Port Elizabeth. The animals were later released at Robben Island and four of them promptly swam back to St. Croix Island, surprising scientists.
In 1983, the exposure of penguins of Dassen Island to the oil slick from the Castillo de Bellver was also a topic of concern given the penguins' conservation status at the time, but owing to the prevailing wind and current, only gannets were oiled.
1994 MV Apollo Sea disaster
African penguin casualties were significant following the sinking of the MV Apollo Sea and a subsequent oil slick in 1994. 10,000 penguins were collected and cleaned, of which less than half survived.
2000 MV Treasure crisis
Disaster struck on 23 June 2000, when the iron ore tanker MV Treasure sank between Robben Island and Dassen Island, South Africa. It released of fuel oil, causing an unprecedented coastal bird crisis and oiling 19,000 adult penguins at the height of the best breeding season on record for this vulnerable species. The oiled birds were brought to an abandoned train repair warehouse in Cape Town to be cared for. An additional 19,500 un-oiled penguins were removed from Dassen Island and other areas before they became oiled and were released about 800 kilometres east of Cape Town. This gave workers enough time to clean up the oiled waters and shores before the birds could complete their long swim home (which took the penguins between one and three weeks). Some of the penguins were named and radio-tracked as they swam back to their breeding grounds. Tens of thousands of volunteers helped with the rescue and rehabilitation process, which was overseen by the International Fund for Animal Welfare (IFAW) and the South African Foundation for the Conservation of Coastal Birds (SANCCOB) and took more than three months to complete. This was the largest animal rescue event in history; more than 91% of the penguins were successfully rehabilitated and released – an amazing feat that could not have been accomplished without such a tremendous international response.
Due to the positive outcome of African penguins being raised in captivity after tragedies such as the Treasure oil spill, the species is considered a good "candidate for a captive-breeding programme which aims to release offspring into the wild"; however, worry about the spread of new strains of avian malaria is a major concern in the situation.
Bringing the birds inland led to the exposure of penguins to parasites and disease vectors such as mosquitoes carrying avian malaria, which has caused 27% of the rehabilitated penguin deaths annually.
2016 & 2019 Port of Ngqura
Small-scale oil spills (of less than ) have occurred at the Port of Ngqura since bunkering activities started there in 2016. Bunkering is a ship refuelling process that can result in oil spills and oil slicks entering the water. Hundreds of African penguins have been harmed following these spills due to the port's close proximity to penguin rookeries on St. Croix Island and seabird habitat on neighbouring Jahleel and Brenton Islands.
Competition with fisheries
Commercial fisheries of sardines and anchovy, the two main prey species of the penguins, have forced these penguins to search for prey farther offshore, as well as having to switch to eating less nutritious prey. Restricting commercial fishing near colony sites such as Robben Island for short periods (3 years) was shown to markedly improve penguin breeding success. Longer closure periods and closures near other colonies are being evaluated.
Conservation status
The African penguin is one of the species to which the African-Eurasian Waterbird Agreement (AEWA) applies. In September 2010, it was listed as endangered under the US Endangered Species Act. As of 2024, the African penguin is listed as critically endangered on the IUCN Red List, with the remaining mature individuals around 19,800 birds in a declining population.
Mediation efforts
Many organisations such as SANCCOB, Dyer Island Conservation Trust, SAMREC, The National Aviary in Pittsburgh, and Raggy Charters with the Penguin Research Fund in Port Elizabeth are working to halt the decline of the African penguin. Measures include: monitoring population trends, hand-rearing and releasing abandoned chicks, establishing artificial nests and proclaiming marine reserves in which fishing is prohibited. Some colonies (such as on Dyer Island) are suspected to be under heavy pressure from predation by Cape fur seals and may benefit from the culling of individual problem animals, which has been found effective (although requiring a large amount of management effort) in trials.
Established in 1968, SANCCOB is currently the only organisation mandated by the South African government to respond to crises involving seabirds along South Africa's coastline and is internationally recognised for the role it played during the MV Treasure oil spill. A modelling exercise conducted in 2003 by the University of Cape Town's FitzPatrick Institute of African Ornithology found that rehabilitating oiled African penguins has resulted in the current population being 19% larger than it would have been in the absence of SANCCOB's rehabilitation efforts.
In February 2015, the Dyer Island Conservation Trust opened the African Penguin and Seabird Sanctuary (APSS) in Gansbaai, South Africa. The centre was opened by then-Department of Tourism minister Derek Hanekom and will serve as a hub for seabird research carried out by the Dyer Island Conservation Trust. The centre will also run local education projects, host international marine volunteers and seek to improve seabird handling techniques and rehabilitation protocols.
Captivity
African penguins are a commonly seen species in zoos across the world. Because they do not require particularly low temperatures, they are often kept in outside enclosures. They adapt fairly well to this captive environment and are rather easy to breed compared to other species of the family. In Europe, the breeding programme EAZA is regulated by Artis Royal Zoo in the Netherlands, whilst in the United States the SSP programme is cooperatively managed by the AZA. The idea is to create a backup captive population, as well as to aid in the conservation of the population in its natural habitat. Between 2010 and 2013, American zoos spent $300,000 on in situ (wild population) conservation.
| Biology and health sciences | Sphenisciformes | Animals |
318262 | https://en.wikipedia.org/wiki/Pont%20Neuf | Pont Neuf | The Pont Neuf (, "New Bridge") is the oldest standing bridge across the river Seine in Paris, France. It stands by the western (downstream) point of the Île de la Cité, the island in the middle of the river that was, between 250 and 225 BCE, the birthplace of Paris, then known as Lutetia and, during the medieval period, the heart of the city.
The bridge is composed of two separate spans, one of five arches joining the left bank to the Île de la Cité, another of seven joining the island to the right bank. Old engraved maps of Paris show that the newly built bridge just grazed the downstream tip of the Île de la Cité; since then, the natural sandbar building of a mid-river island, aided by stone-faced embankments called quais, has extended the island. Today the tip of the island is the location of the Square du Vert-Galant, a small public park named in honour of Henry IV, nicknamed the "Green Gallant".
The name Pont Neuf was given to distinguish it from older bridges that were lined on both sides with houses, and has remained after all of those were replaced. Its name notwithstanding, it has long been the oldest bridge in Paris crossing the Seine. It has been listed since 1889 as a monument historique by the French Ministry of Culture.
Construction
As early as 1550, Henry II considered building a new bridge at the Ile de la Cite because the existing Pont Notre-Dame was congested and needed repair. The idea was not advanced for lack of funds. By 1577, however, Henry III released funds from the national treasury for a new bridge and appointed a building commission for its designing and planning. Henry rejected the first design proposed by the committee, which included monumental arches, but no plan for buildings along the sides. The commission proceeded in 1578 with modifications to its initial plan, perhaps devised by the royal architect, Androuet de Cerceau. While Henry had already allowed for piers to be driven for the northern arm of bridge, the first construction under the 1579 design indicated a wider deck in preparation of buildings to be constructed on the side. The houses were never built, but the wide bridge deck was retained.
In February 1578, the decision to build the bridge was made by Henry III who laid its first stone in on 31 May 1578, the same year when the foundations of four piers and one abutment were completed. Pierre des Isles, one of the builders, convinced the supervisory commission that the bridge, which was originally planned straight, would be more resistant to the river currents if its two sections were built at a slight angle. The change was adopted in May 1578.
Further design changes were made during the summer of 1579. First, the number of arches was changed from eight and four to seven and five. This was not a problem on the north side, where nothing had been built, but on the south, where the four piles and the abutment on the Left Bank were already laid, the addition of the fifth arch necessitated reducing the length of the platform on the island, the terre-plein, from 28.5 toises to about 19. Second, it was decided to allow houses to be built on the bridge (though they never were). This required the widening of the bridge. The remaining piers were built over the next nine years. After a long delay beginning in 1588, due to political unrest and to the Wars of Religion, construction was resumed in 1599 under the reign of Henry IV. The bridge was opened to traffic in 1604 and completed in July 1606. It was inaugurated by Henry IV in 1607.
Like most bridges of its time, the Pont Neuf is constructed as a series of many short arch bridges, following Roman precedents. It was the first stone bridge in Paris not to support houses in addition to a thoroughfare, and was also fitted with pavements protecting pedestrians from mud and horses; pedestrians could also step aside into its bastions to let a bulky carriage pass. The decision not to include houses on the bridge can be traced back directly to Henry IV, who decided against their inclusion on the grounds that houses would impede a clear view of the Louvre, which the newly built galerie du bord de l'eau linked to the Tuileries Palace.
Pont Neuf was for a long time the widest bridge in Paris. It has undergone much repair and renovation work, including rebuilding of seven spans in the long arm and lowering of the roadway by changing the arches from an almost semi-circular to elliptical form (1848–1855), lowering of sidewalks and faces of the piers, spandrels, cornices and replacing crumbled corbels as closely to the originals as possible. In 1885, one of the piers of the short arm was undermined, removing the two adjacent arches, requiring them to be rebuilt and all the foundations strengthened.
A major restoration of the Pont Neuf was begun in 1994 and was completed in 2007, the year of its 400th anniversary.
Mascarons
The mascarons are the stone masks, 381 in number, each being different and which decorate the sides of the bridge. They represent the heads of forest and field divinities from ancient mythology, as well as satyrs and sylvains. They are copies of the originals attributed to the French Renaissance sculptor Germain Pilon (1525–1590), who also sculpted the tomb of King Henry II of France and Queen Catherine de'Medici in the Basilica of St Denis, five kilometers north of Paris. The mascarons remained in place until 1851–1854, when the bridge was completely rebuilt. At that time six of the original mascarons from the 16th century were placed in the Musée Carnavalet, along with eight molds of other originals. Eight other originals were first placed in the Musée de Cluny – Musée national du Moyen Âge, and are now in the French National Museum of the Renaissance in the Château d'Écouen. During their reconstruction, the Renaissance masks were replaced with copies made by noted 19th-century sculptors, including Hippolyte Maindron, Hubert Lavigne, Antoine-Louis Barye and Fontenelle. Fontenelle made 61 masks, which are found on the upstream side of the bridge between the right bank and the Île de la Cité.
Equestrian statue of Henry IV
At the point where the bridge crosses the Île de la Cité, there stands a bronze equestrian statue of king Henry IV, originally commissioned from Giambologna under the orders of Marie de Médicis, Henri's widow and Regent of France. After his death, Giambologna's assistant Pietro Tacca completed the statue, which was erected on its pedestal by Pietro Francavilla, in 1614. It was destroyed in 1792 during the French Revolution, but was rebuilt in 1818, following the restoration of the Bourbon monarchy. Commissioned from public donations, bronze for the new statue was obtained from a statue of Louis Charles Antoine Desaix and melted down. The new statue was cast from a mold made using a surviving cast of the original. Inside the statue, the new sculptor François-Frédéric Lemot put four boxes, containing a history of the life of Henry IV, a 17th-century parchment certifying the original statue, a document describing how the new statue was commissioned, and a list of people who contributed to a public subscription.
La Samaritaine
Between 1712 and 1719, replacing an earlier one, a large pump house was built on the bridge. It was decorated with an image of the Samaritan woman at the well. As a result, the structure (which included a carillon) was named La Samaritaine. Years after it was torn down (in 1813), Ernest Cognacq, a 19th-century merchant, set up a stand on the site and gradually grew his business to what became, in 1869, the department store La Samaritaine.
As the centre of Paris
Upon completion, Pont Neuf attracted throngs of visitors, many of whom used the bridge as a public square, conducting business, socializing, and taking in the view. One contemporary writer repeated a proverb about Pont Neuf to illustrate the variety of people who frequented the bridge, "one never crossed the Pont Neuf without meeting three things: a monk, a girl and a white horse."
All through the 18th century, the Pont Neuf was the center of Paris, lively with both crime and commerce: Czar Peter the Great, who came to study French civilization under the regency of the Duke d'Orleans, declared that he had found nothing more curious in Paris than the Pont Neuf; and, sixty years later, the philosopher Franklin wrote to his friends in America that he had not understood the Parisian character except in crossing the Pont Neuf.
In 1862, Édouard Fournier traced its history in his lively two-volume Histoire du Pont-Neuf. He describes how, even before it was completed (in 1607), gangs hid out in and around it, robbing and murdering people. It remained a dangerous place even as it became busier. For a long time, the bridge even had its own gallows.
This did not prevent people from congregating there, drawn by various stands and street performers (acrobats, fire-eaters, musicians, etc.). Charlatans and quacks of various sorts were also common, as well as the hustlers (shell game hucksters, etc.) and pickpockets often found in crowds – not to mention a lively trade in prostitution. Among the many businesses which, however, unofficially set up there, were several famous tooth pullers.
In 1701, Cotolendi quoted a letter supposedly written by a Sicilian tourist:
One finds on the Pont-Neuf an infinity of people who give tickets, some put fallen teeth back in, and others make crystal eyes; there are those who cure incurable illnesses; those who claim to have discovered the virtues of some powdered stones to white and to beautify the face. This one claims he makes old men young; there are those who remove wrinkles from the forehead and the eyes, who make wooden legs to repair the violence of bombs; finally everybody is so applied to work, so strongly and continually, that the devil can tempt no one but on Holidays and Sundays.
With its numerous sellers of pamphlets and satirical performers, it was also a center for social commentary:
In the 16th cent. the Pont-Neuf was the scene of the recitals of Tabarin, a famous satirist of the day, and it was long afterwards the favourite rendezvous of news-vendors, jugglers, showmen, loungers, and thieves. Any popular witticism in verse was long known as un Pont-Neuf.
In the seventeenth century, that bridge of memories, the old Pont Neuf of Paris, was the rendezvous of quacksalvers and mountebanks. Booths for the sale of various articles lined the sides of the bridge. People flocked there to see the sights, laugh, chat, make love and enjoy life as only Parisians can. Students and grisettes of the Quartier latin elbowed ladies and gentlemen of the court. Bourgeois families came to study the flippant manners of the aristocrats. Poodle clippers plied their trade; jugglers amused the quid nuncs with feats of dexterity; traveling dentists pulled teeth and sold balsams; clowns tumbled; and last, but not least, pickpockets lifted purses and silk handkerchiefs with impunity. Says Augustus J. C. Hare (Walks in Paris): "So central an artery is the Pont Neuf, that it used to be a saying with the Parisian police, that if, after watching three days, they did not see a man cross the bridge, he must have left Paris." One of the principal vendors of quack nostrums of the Pont Neuf was Montdor. He was aided by a buffoon named Tabarin, who made facetious replies to questions asked by his master, accompanied with laughable grimaces and grotesque gestures. The modern ringmaster and clown of the circus have similar scenes together, minus the selling of medicines.
Under Louis XV, thieves and entertainers were joined by recruiters, or "sellers of human flesh", who did their best to lure newcomers to Paris and others "with as much violence as the sale of Negros in the Congo". Silversmiths and other luxury businesses nearby (which gave their name to the Quai des Orfèvres) drew visitors as well.
One yearly event, held on the nearby Place Dauphine, prefigured the Salon des Refusés which would give rise to the Impressionists. During the celebration of the Corpus Christi (Fête-Dieu), the Place Dauphine hosted one of the most magnificent reposoirs (portable altars for the Host).
Along with all the rich silverwork and tapestries placed on it, some local silversmiths ordered paintings for these. This led to art dealers being asked to participate and, ultimately, to the newest talents being shown at the Petite Fête-Dieu (the Small Corpus Christi), a reduced version of the Corpus Christi holiday which took place eight days later. Though their canvases were only shown from six in the morning to noon, this became an important opportunity for unknown artists to draw attention. Among other things, this led to the painters there signing their work, as was not frequent in the Salon – which was not always an advantage when the work was publicly and loudly critiqued.
Showing works, which often had no pretense of a religious subject, they might then be noticed and find an entree into the official Academy. Chardin is one of the most famous painters to have started this way.
In 1720, a young man of about twenty-two, son of the man who maintained the king's billiards, displayed a canvas here showing an antique bas-relief. J.-B. Vanloo passed by, looked at the canvas for a long time, found great qualities there, and bought it. He wanted afterwards to know the young painter, encouraged him, gave him advice, of which the latter perhaps had no need, got him work, which was more useful, and eight years later, the unknown of the place Dauphine was his colleague at the Academy of Painting.... he was called Jean-Baptiste-Siméon Chardin.
The slow decline of the bridge's central role began in 1754: "Starting in 1754, the first year of the vogue, the madness of the boulevards, it was no longer the thing to talk about the Cours [the Champs-Elysées], and still less of this poor Pont-Neuf. To the Boulevard, at once, long live the Boulevard!". Still the bridge remained a lively place through the end of the century. With time, people became wary of its reputation and other changes subdued its atmosphere. In 1840, Lacroix wrote: "Once the pont Neuf was a perpetual fair; at present, it is just a bridge to be crossed without stopping."
Possible first photograph of human being
In 1838, Louis Daguerre produced his famous daguerreotype portrait of the View of the Boulevard du Temple, widely considered the first photograph where a human can be seen. However, between 1836 and 1837, Daguerre made several tests, in order to experiment with and perfect the new technique in an outdoor environment.
One surviving example is an image of the Pont Neuf and the equestrian statue of Henry IV, made possibly as early as 1836. On the lower-left side of the image, what appears to be a worker, or perhaps two, can be seen lying against the fence, in the shadow of the statue.
Christo's project
In 1985, after years of negotiation with the mayor of Paris, the art duo Christo and Jeanne-Claude wrapped the Pont Neuf.
Access
| Technology | Bridges | null |
318265 | https://en.wikipedia.org/wiki/Moschidae | Moschidae | Moschidae is a family of pecoran even-toed ungulates, containing the musk deer (Moschus) and its extinct relatives. They are characterized by long "saber teeth" instead of horns, antlers or ossicones, modest size (Moschus only reaches ; other taxa were even smaller) and a lack of facial glands. While various Oligocene and Miocene pecorans were previously assigned to this family, recent studies find that most should be assigned to their own clades, although further research would need to confirm these traits. As a result, Micromeryx, Hispanomeryx, and Moschus are the only undisputed moschid members, making them known from at least 18 Ma. The group was abundant across Eurasia and North America during the Miocene, but afterwards declined to only the extant genus Moschus by the early Pleistocene.
Taxonomy and classification
Until the early 21st century, it was believed that the musk deer (family Moschidae) were an adjacent, sister-group to the "true deer" of the family Cervidae (caribou, moose, elk, and roughly 40–50 other species); however, a 2003 phylogenetic study by Alexandre Hassanin (of the National Museum of Natural History, France) and co., based on mitochondrial and nuclear analyses, revealed that Moschidae and Bovidae (antelope, cattle, goats, sheep), together, form a sister-clade to Cervidae. According to the study, the Cervidae diverged from the Bovidae-Moschidae clade roughly 27–28 million years ago. The following cladogram is based on this 2003 study:
After Prothero (2007)
Family Moschidae
| Biology and health sciences | Other artiodactyla | Animals |
318270 | https://en.wikipedia.org/wiki/Musk%20deer | Musk deer | Musk deer can refer to any one, or all eight, of the species that make up Moschus, the only extant genus of the family Moschidae. Despite being commonly called deer, they are not true deer belonging to the family Cervidae, but rather their family is closely related to Bovidae, the group that contains antelopes, bovines, sheep, and goats. The musk deer family differs from cervids, or true deer, by lacking antlers and preorbital glands also, possessing only a single pair of teats, a gallbladder, a caudal gland, a pair of canine tusks and—of particular economic importance to humans—a musk gland.
Musk deer live mainly in forested and alpine scrub habitats in the mountains of South Asia, notably the Himalayas. Moschids, the proper term when referring to this type of deer rather than one/multiple species of musk deer, are entirely Asian in their present distribution, being extinct in Europe where the earliest musk deer are known to have existed from Oligocene deposits.
Characteristics
Musk deer resemble small deer, with a stocky build and hind legs longer than their front legs. They are about long, high at the shoulder, and weigh between . The feet of musk deer are adapted for climbing in rough terrain. Like the Chinese water deer, a cervid, they have no antlers, but the males do have enlarged upper canines, forming sabre-like tusks. The dental formula is similar to that of true deer:
The musk gland is found only in adult males. It lies in a sac located between the genitals and the umbilicus, and its secretions are most likely used to attract mates.
Musk deer are herbivores, living in hilly, forested environments, generally far from human habitation. Like true deer, they eat mainly leaves, flowers, and grasses, with some mosses and lichens. They are solitary animals and maintain well-defined territories, which they scent mark with their caudal glands. Musk deer are generally shy and either nocturnal or crepuscular.
Males leave their territories during the rutting season and compete for mates, using their tusks as weapons. In order to indicate their area, musk deer build latrines. These locations can be used to identify the musk deer's existence, number, and preferred habitat in the wild. Female musk deer give birth to a single fawn after about 150–180 days. The newborn young are very small and essentially motionless for the first month of their lives, a feature that helps them remain hidden from predators.
Musk deer have been hunted for their scent glands, which are used in perfumes. The glands can fetch up to $45,000/kg on the black market. It is rumored that ancient royalty wore the scent of the musk deer, and that it is an aphrodisiac.
Population
Musk deer have a global population between 400,000 to 800,000 currently, however the exact count is undetermined. They are widely spread; however, their population density increases within China, Russia, and Mongolia. Musk deer are commonly found in China, and they are spread over 17 provinces. This population is mainly located around the Himalayas in southern Asia, southeast Asia, and eastern Asia. They are also found in a few spots in Russia. As of 2003, they became a protected species due to their declined overall population. Musk deer have many subspecies that have varying population sizes, within the overall total, and all are threatened. Over the past twenty years, the populations have been able to slightly recover due to the captive breeding of these animals, specifically in China. Musk deer populations are recovering due to the protocols and rules being set in place to protect the species.
Habitat
The musk deer species is generally solitary and lives in the higher regions of mountain ranges, such as the Himalayas. The varying species' habitats include different atmospheres and necessary resources for their survival, while including similar universal resources. Musk deer population has been declining recently due to environmental and human factors. As a large-bodied mammal, they have great needs that are not able to be sustained due to habitat fragmentation. This species is largely protected due to the threat of extinction, due to the increase in illegal hunting. Illegal hunting has significantly decreased the population throughout many of the provinces musk deer occupy. Their habitats are being lost to colonization and deforestation and hunting for musk deer was on the rise. They were hunted for their distinct products that are very valuable in the market. Since then, the Chinese government has stepped in to regulate these issues. They have placed rules pertaining to the killing of musk deer and created havens for the deer to survive. To help with the declining numbers, the deforestation of their natural habitat should be stopped and new habitats should be invested in them. Global climate change has also driven the musk deer population down. The warmer climates result in the drive to higher elevations and latitudes. Global warming and habitat fragmentation are two causes for the population decrease.
Evolution
Musk deer are the only surviving members of the Moschidae, a family with a fossil record extending over 25 million years to the late Oligocene. The group was abundant across Eurasia and North America until the late Miocene, but underwent a substantial decline, with no Pliocene fossil record and Moschus the only genus since the Pleistocene. The oldest records of the genus Moschus are known from the Late Miocene (Turolian) of Lufeng, China.
Taxonomy
While they have been traditionally classified as members of the deer family (as the subfamily "Moschinae") and all the species were classified as one species (under Moschus moschiferus), recent studies have indicated that moschids are more closely related to bovids (antelope, goats, sheep and cattle).
| Biology and health sciences | Other artiodactyla | Animals |
318319 | https://en.wikipedia.org/wiki/Chevrotain | Chevrotain | Chevrotains, or mouse-deer, are diminutive, even-toed ungulates that make up the family Tragulidae, and are the only living members of the infraorder Tragulina. The 10 extant species are placed in three genera, but several species also are known only from fossils. The extant species are found in forests in South and Southeast Asia; a single species, the water chevrotain, is found in the rainforests of Central and West Africa. In November 2019, conservation scientists announced that they had photographed silver-backed chevrotains (Tragulus versicolor) in a Vietnamese forest for the first time since the last confirmed sightings in 1990.
They are solitary, or live in loose groupings or pairs, and feed almost exclusively on plant material. Chevrotains are the smallest hoofed mammals in the world. The Asian species weigh between , while the African chevrotain is considerably larger, at . With an average length of and an average height of , the Java mouse-deer is the smallest surviving ungulate (hoofed) mammal, as well as the smallest artiodactyl (even-toed ungulate). Despite their common name of "mouse deer", they are not closely related to true deer.
Etymology
The word "chevrotain" comes from the Middle French word chevrot (kid or fawn), derived from chèvre (goat).
The single African species is consistently known as "chevrotain". The names "chevrotain" and "mouse-deer" have been used interchangeably among the Asian species, though recent authorities typically have preferred chevrotain for the species in the genus Moschiola and mouse-deer for the species in the genus Tragulus. Consequently, all species with pale-spotted or -striped upper parts are known as "chevrotain" and without are known as "mouse-deer".
The Telugu name for the Indian spotted chevrotain is jarini pandi, which literally means "a deer and a pig". In Kannada, it is called barka (ಬರ್ಕ), in Malayalam, it is called kūramān, and the Konkani name for it is barinka. The Tamil term is sarukumāṉ "leaf-pile deer". The Sinhala name roughly translates to "mouse-like deer". This was used in the scientific name of the Sri Lankan spotted chevrotain, M. meminna.
Biology
The family was widespread and successful from the Oligocene (34 million years ago) through the Miocene (about 5 million years ago), but has remained almost unchanged over that time and remains as an example of an archaic ruminant type. They have four-chambered stomachs to ferment tough plant foods, but the third chamber is poorly developed. Unlike other artiodactyls, they lack an carotid rete, and so cannot heat exchange cool blood entering their brains, a thermoregulatory innovation that allows other artiodactyls to exploit hot arid habitats. Though most species feed exclusively on plant material, the water chevrotain occasionally takes insects and crabs or scavenges meat and fish. Like other ruminants, they lack upper incisors. They give birth to only a single young.
In other respects, however, they have primitive features, closer to nonruminants such as pigs. All species in the family lack antlers and horns, but both sexes have elongated canine teeth. These are especially prominent in males, where they project out on either side of the lower jaw, and are used in fights. Their legs are short and thin, which leave them lacking in agility, but also helps to maintain a smaller profile to aid in running through the dense foliage of their environments. Other pig-like features include the presence of four toes on each foot, the absence of facial scent glands, premolars with sharp crowns, and the form of their sexual behaviour and copulation.
They are solitary or live in pairs. The young are weaned at three months of age, and reach sexual maturity between 5 and 10 months, depending on species. Parental care is relatively limited. Although they lack the types of scent glands found in most other ruminants, they do possess a chin gland for marking each other as mates or antagonists, and, in the case of the water chevrotain, anal and preputial glands for marking territory. Their territories are relatively small, on the order of , but neighbors generally ignore each other, rather than compete aggressively.
Some of the species show a remarkable affinity with water, often remaining submerged for prolonged periods to evade predators or other unwelcome intrusions. This has also lent support to the idea that whales evolved from water-loving creatures that looked like small deer.
Taxonomy
Tragulidae's placement within Artiodactyla can be represented in the following cladogram:
Traditionally, only four extant species were recognized in the family Tragulidae. In 2004, T. nigricans and T. versicolor were split from T. napu, and T. kanchil and T. williamsoni were split from T. javanicus. In 2005, M. indica and M. kathygre were split from M. meminna. With these changes, the 10 extant species are:
Family Tragulidae
Genus Hyemoschus
Water chevrotain, Hyemoschus aquaticus
Genus Moschiola
Indian spotted chevrotain, Moschiola indica
Sri Lankan spotted chevrotain, Moschiola meminna
Yellow-striped chevrotain, Moschiola kathygre
Genus Tragulus
Java mouse-deer, Tragulus javanicus
Lesser mouse-deer or kanchil, Tragulus kanchil
Greater mouse-deer, Tragulus napu
Philippine mouse-deer, Tragulus nigricans
Vietnam mouse-deer, Tragulus versicolor
Williamson's mouse-deer, Tragulus williamsoni
Ancient chevrotains
The Hypertragulidae were closely related to the Tragulidae.
The six extinct chevrotain genera include:
Genus Dorcatherium
Dorcatherium minus from Pakistan
Dorcatherium majus from Pakistan
Dorcatherium naui, from Central Europe
Genus Dorcabune
Dorcabune anthracotherioides from Pakistan
Dorcabune nagrii from Pakistan
Genus Afrotragulus Sánchez, Quiralte, Morales and Pickford, 2010
Afrotragulus moruorotensis (previously "Dorcatherium" moruorotensis Pickford, 2001) (early Miocene) from Moruorot, Kenya
Afrotragulus parvus (previously "D." parvus Withworth 1958) (early Miocene) from Rusinga Island, Kenya
Genus Siamotragulus
Siamotragulus sanyathanai Thomas, Ginsburg, Hintong and Suteethorn, 1990 (middle Miocene) from Lampang, Thailand
Siamotragulus haripounchai Mein and Ginsburg, 1997 (Miocene) from Lamphun, Thailand
Genus Yunnanotherium
Genus Archaeotragulus
Archaeotragulus krabiensis Metais, Chaimanee, Jaeger and Ducrocq, 2001 (late Eocene) from Krabi, Thailand
The extinct chevrotains might also include
Genus Krabitherium
Krabitherium waileki Metais, Chaimanee, Jaeger and Ducrocq, 2007 (late Eocene) from Krabi, Thailand
Genus Nalameryx
Nalameryx savagei
Nalameryx sulaimani
Mythology
According to the Malay Annals, King Parameswara, seeking a place to found a new city, came to a place where he saw a mouse deer (kancil in Malay) kicking his hunting dog into the river. He thought this boded well, remarking, 'this place is excellent, even the mouse deer is formidable; it is best that we establish a kingdom here'. He then founded there the city of Malacca. In memory of this founding legend, the coat of arms of Malacca depicts two mouse deer.
The mouse deer or Sang Kancil is also a clever character from several Malay folktales.
| Biology and health sciences | Other artiodactyla | Animals |
318378 | https://en.wikipedia.org/wiki/Crane%20%28machine%29 | Crane (machine) | A crane is a machine used to move materials both vertically and horizontally, utilizing a system of a boom, hoist, wire ropes or chains, and sheaves for lifting and relocating heavy objects within the swing of its boom. The device uses one or more simple machines, such as the lever and pulley, to create mechanical advantage to do its work. Cranes are commonly employed in transportation for the loading and unloading of freight, in construction for the movement of materials, and in manufacturing for the assembling of heavy equipment.
The first known crane machine was the shaduf, a water-lifting device that was invented in ancient Mesopotamia (modern Iraq) and then appeared in ancient Egyptian technology. Construction cranes later appeared in ancient Greece, where they were powered by men or animals (such as donkeys), and used for the construction of buildings. Larger cranes were later developed in the Roman Empire, employing the use of human treadwheels, permitting the lifting of heavier weights. In the High Middle Ages, harbour cranes were introduced to load and unload ships and assist with their construction—some were built into stone towers for extra strength and stability. The earliest cranes were constructed from wood, but cast iron, iron and steel took over with the coming of the Industrial Revolution.
For many centuries, power was supplied by the physical exertion of men or animals, although hoists in watermills and windmills could be driven by the harnessed natural power. The first mechanical power was provided by steam engines, the earliest steam crane being introduced in the 18th or 19th century, with many remaining in use well into the late 20th century. Modern cranes usually use internal combustion engines or electric motors and hydraulic systems to provide a much greater lifting capability than was previously possible, although manual cranes are still utilized where the provision of power would be uneconomic.
There are many different types of cranes, each tailored to a specific use. Sizes range from the smallest jib cranes, used inside workshops, to the tallest tower cranes, used for constructing high buildings. Mini-cranes are also used for constructing high buildings, to facilitate constructions by reaching tight spaces. Large floating cranes are generally used to build oil rigs and salvage sunken ships.
Some lifting machines do not strictly fit the above definition of a crane, but are generally known as cranes, such as stacker cranes and loader cranes.
Etymology
Cranes were so called from the resemblance to the long neck of the bird, cf. , French grue.
History
Ancient Near East
The first type of crane machine was the shadouf, which had a lever mechanism and was used to lift water for irrigation. It was invented in Mesopotamia (modern Iraq) circa 3000 BC. The shadouf subsequently appeared in ancient Egyptian technology circa 2000 BC.
Ancient Greece
A crane for lifting heavy loads was developed by the Ancient Greeks in the late 6th century BC. The archaeological record shows that no later than c. 515 BC distinctive cuttings for both lifting tongs and lewis irons begin to appear on stone blocks of Greek temples. Since these holes point at the use of a lifting device, and since they are to be found either above the center of gravity of the block, or in pairs equidistant from a point over the center of gravity, they are regarded by archaeologists as the positive evidence required for the existence of the crane.
The introduction of the winch and pulley hoist soon led to a widespread replacement of ramps as the main means of vertical motion. For the next 200 years, Greek building sites witnessed a sharp reduction in the weights handled, as the new lifting technique made the use of several smaller stones more practical than fewer larger ones. In contrast to the archaic period with its pattern of ever-increasing block sizes, Greek temples of the classical age like the Parthenon invariably featured stone blocks weighing less than 15–20 metric tons. Also, the practice of erecting large monolithic columns was practically abandoned in favour of using several column drums.
Although the exact circumstances of the shift from the ramp to the crane technology remain unclear, it has been argued that the volatile social and political conditions of Greece were more suitable to the employment of small, professional construction teams than of large bodies of unskilled labour, making the crane preferable to the Greek polis over the more labour-intensive ramp which had been the norm in the autocratic societies of Egypt or Assyria.
The first unequivocal literary evidence for the existence of the compound pulley system appears in the Mechanical Problems (Mech. 18, 853a32–853b13) attributed to Aristotle (384–322 BC), but perhaps composed at a slightly later date. Around the same time, block sizes at Greek temples began to match their archaic predecessors again, indicating that the more sophisticated compound pulley must have found its way to Greek construction sites by then.
Roman Empire
The heyday of the crane in ancient times came during the Roman Empire, when construction activity soared and buildings reached enormous dimensions. The Romans adopted the Greek crane and developed it further. There is much available information about their lifting techniques, thanks to rather lengthy accounts by the engineers Vitruvius (De Architectura 10.2, 1–10) and Heron of Alexandria (Mechanica 3.2–5). There are also two surviving reliefs of Roman treadwheel cranes, with the Haterii tombstone from the late first century AD being particularly detailed.
The simplest Roman crane, the trispastos, consisted of a single-beam jib, a winch, a rope, and a block containing three pulleys. Having thus a mechanical advantage of 3:1, it has been calculated that a single man working the winch could raise (3 pulleys x = 150), assuming that represent the maximum effort a man can exert over a longer time period. Heavier crane types featured five pulleys (pentaspastos) or, in case of the largest one, a set of three by five pulleys (Polyspastos) and came with two, three or four masts, depending on the maximum load. The polyspastos, when worked by four men at both sides of the winch, could readily lift (3 ropes x 5 pulleys x 4 men x = ). If the winch was replaced by a treadwheel, the maximum load could be doubled to at only half the crew, since the treadwheel possesses a much bigger mechanical advantage due to its larger diameter. This meant that, in comparison to the construction of the ancient Egyptian pyramids, where about 50 men were needed to move a 2.5 ton stone block up the ramp ( per person), the lifting capability of the Roman polyspastos proved to be 60 times higher ( per person).
However, numerous extant Roman buildings which feature much heavier stone blocks than those handled by the polyspastos indicate that the overall lifting capability of the Romans went far beyond that of any single crane. At the temple of Jupiter at Baalbek, for instance, the architrave blocks weigh up to 60 tons each, and one corner cornice block even over 100 tons, all of them raised to a height of about . In Rome, the capital block of Trajan's Column weighs 53.3 tons, which had to be lifted to a height of about (see construction of Trajan's Column).
It is assumed that Roman engineers lifted these extraordinary weights by two measures (see picture below for comparable Renaissance technique): First, as suggested by Heron, a lifting tower was set up, whose four masts were arranged in the shape of a quadrangle with parallel sides, not unlike a siege tower, but with the column in the middle of the structure (Mechanica 3.5). Second, a multitude of capstans were placed on the ground around the tower, for, although having a lower leverage ratio than treadwheels, capstans could be set up in higher numbers and run by more men (and, moreover, by draught animals). This use of multiple capstans is also described by Ammianus Marcellinus (17.4.15) in connection with the lifting of the Lateranense obelisk in the Circus Maximus (c. 357 AD). The maximum lifting capability of a single capstan can be established by the number of lewis iron holes bored into the monolith. In case of the Baalbek architrave blocks, which weigh between 55 and 60 tons, eight extant holes suggest an allowance of 7.5 ton per lewis iron, that is per capstan. Lifting such heavy weights in a concerted action required a great amount of coordination between the work groups applying the force to the capstans.
Middle Ages
During the High Middle Ages, the treadwheel crane was reintroduced on a large scale after the technology had fallen into disuse in western Europe with the demise of the Western Roman Empire. The earliest reference to a treadwheel (magna rota) reappears in archival literature in France about 1225, followed by an illuminated depiction in a manuscript of probably also French origin dating to 1240. In navigation, the earliest uses of harbor cranes are documented for Utrecht in 1244, Antwerp in 1263, Bruges in 1288 and Hamburg in 1291, while in England the treadwheel is not recorded before 1331.
Generally, vertical transport could be done more safely and inexpensively by cranes than by customary methods. Typical areas of application were harbors, mines, and, in particular, building sites where the treadwheel crane played a pivotal role in the construction of the lofty Gothic cathedrals. Nevertheless, both archival and pictorial sources of the time suggest that newly introduced machines like treadwheels or wheelbarrows did not completely replace more labor-intensive methods like ladders, hods and handbarrows. Rather, old and new machinery continued to coexist on medieval construction sites and harbors.
Apart from treadwheels, medieval depictions also show cranes to be powered manually by windlasses with radiating spokes, cranks and by the 15th century also by windlasses shaped like a ship's wheel. To smooth out irregularities of impulse and get over 'dead-spots' in the lifting process flywheels are known to be in use as early as 1123.
The exact process by which the treadwheel crane was reintroduced is not recorded, although its return to construction sites has undoubtedly to be viewed in close connection with the simultaneous rise of Gothic architecture. The reappearance of the treadwheel crane may have resulted from a technological development of the windlass from which the treadwheel structurally and mechanically evolved. Alternatively, the medieval treadwheel may represent a deliberate reinvention of its Roman counterpart drawn from Vitruvius' De architectura which was available in many monastic libraries. Its reintroduction may have been inspired, as well, by the observation of the labor-saving qualities of the waterwheel with which early treadwheels shared many structural similarities.
Structure and placement
The medieval treadwheel was a large wooden wheel turning around a central shaft with a treadway wide enough for two workers walking side by side. While the earlier 'compass-arm' wheel had spokes directly driven into the central shaft, the more advanced "clasp-arm" type featured arms arranged as chords to the wheel rim, giving the possibility of using a thinner shaft and providing thus a greater mechanical advantage.
Contrary to a popularly held belief, cranes on medieval building sites were neither placed on the extremely lightweight scaffolding used at the time nor on the thin walls of the Gothic churches which were incapable of supporting the weight of both hoisting machine and load. Rather, cranes were placed in the initial stages of construction on the ground, often within the building. When a new floor was completed, and massive tie beams of the roof connected the walls, the crane was dismantled and reassembled on the roof beams from where it was moved from bay to bay during construction of the vaults. Thus, the crane "grew" and "wandered" with the building with the result that today all extant construction cranes in England are found in church towers above the vaulting and below the roof, where they remained after building construction for bringing material for repairs aloft.
Less frequently, medieval illuminations also show cranes mounted on the outside of walls with the stand of the machine secured to putlogs.
Mechanics and operation
In contrast to modern cranes, medieval cranes and hoists — much like their counterparts in Greece and Rome — were primarily capable of a vertical lift, and not used to move loads for a considerable distance horizontally as well. Accordingly, lifting work was organized at the workplace in a different way than today. In building construction, for example, it is assumed that the crane lifted the stone blocks either from the bottom directly into place, or from a place opposite the centre of the wall from where it could deliver the blocks for two teams working at each end of the wall. Additionally, the crane master who usually gave orders at the treadwheel workers from outside the crane was able to manipulate the movement laterally by a small rope attached to the load. Slewing cranes which allowed a rotation of the load and were thus particularly suited for dockside work appeared as early as 1340. While ashlar blocks were directly lifted by sling, lewis or devil's clamp (German Teufelskralle), other objects were placed before in containers like pallets, baskets, wooden boxes or barrels.
It is noteworthy that medieval cranes rarely featured ratchets or brakes to forestall the load from running backward. This curious absence is explained by the high friction force exercised by medieval tread-wheels which normally prevented the wheel from accelerating beyond control.
Harbour usage
According to the "present state of knowledge" unknown in antiquity, stationary harbor cranes are considered a new development of the Middle Ages. The typical harbor crane was a pivoting structure equipped with double treadwheels. These cranes were placed docksides for the loading and unloading of cargo where they replaced or complemented older lifting methods like see-saws, winches and yards.
Two different types of harbor cranes can be identified with a varying geographical distribution: While gantry cranes, which pivoted on a central vertical axle, were commonly found at the Flemish and Dutch coastside, German sea and inland harbors typically featured tower cranes where the windlass and treadwheels were situated in a solid tower with only jib arm and roof rotating. Dockside cranes were not adopted in the Mediterranean region and the highly developed Italian ports where authorities continued to rely on the more labor-intensive method of unloading goods by ramps beyond the Middle Ages.
Unlike construction cranes where the work speed was determined by the relatively slow progress of the masons, harbor cranes usually featured double treadwheels to speed up loading. The two treadwheels whose diameter is estimated to be 4 m or larger were attached to each side of the axle and rotated together. Their capacity was 2–3 tons, which apparently corresponded to the customary size of marine cargo. Today, according to one survey, fifteen treadwheel harbor cranes from pre-industrial times are still extant throughout Europe. Some harbour cranes were specialised at mounting masts to newly built sailing ships, such as in Gdańsk, Cologne and Bremen. Beside these stationary cranes, floating cranes, which could be flexibly deployed in the whole port basin came into use by the 14th century.
A sheer hulk (or shear hulk) was used in shipbuilding and repair as a floating crane in the days of sailing ships, primarily to place the lower masts of a ship under construction or repair. Booms known as sheers were attached to the base of a hulk's lower masts or beam, supported from the top of those masts. Blocks and tackle were then used in such tasks as placing or removing the lower masts of the vessel under construction or repair. These lower masts were the largest and most massive single timbers aboard a ship, and erecting them without the assistance of either a sheer hulk or land-based masting sheer was extremely difficult.
The concept of sheer hulks originated with the Royal Navy in the 1690s, and persisted in Britain until the early nineteenth century. Most sheer hulks were decommissioned warships; Chatham, built in 1694, was the first of only three purpose-built vessels. There were at least six sheer hulks in service in Britain at any time throughout the 1700s. The concept spread to France in the 1740s with the commissioning of a sheer hulk at the port of Rochefort.
Early modern age
A lifting tower similar to that of the ancient Romans was used to great effect by the Renaissance architect Domenico Fontana in 1586 to relocate the 361 t heavy Vatican obelisk in Rome. From his report, it becomes obvious that the coordination of the lift between the various pulling teams required a considerable amount of concentration and discipline, since, if the force was not applied evenly, the excessive stress on the ropes would make them rupture.
Cranes were also used domestically during this period. The chimney or fireplace crane was used to swing pots and kettles over the fire and the height was adjusted by a trammel.
Industrial revolution
With the onset of the Industrial Revolution the first modern cranes were installed at harbours for loading cargo. In 1838, the industrialist and businessman William Armstrong designed a water-powered hydraulic crane. His design used a ram in a closed cylinder that was forced down by a pressurized fluid entering the cylinder and a valve regulated the amount of fluid intake relative to the load on the crane. This mechanism, the hydraulic jigger, then pulled on a chain to lift the load.
In 1845 a scheme was set in motion to provide piped water from distant reservoirs to the households of Newcastle. Armstrong was involved in this scheme and he proposed to Newcastle Corporation that the excess water pressure in the lower part of town could be used to power one of his hydraulic cranes for the loading of coal onto barges at the Quayside. He claimed that his invention would do the job faster and more cheaply than conventional cranes. The corporation agreed to his suggestion, and the experiment proved so successful that three more hydraulic cranes were installed on the Quayside.
The success of his hydraulic crane led Armstrong to establish the Elswick works at Newcastle, to produce his hydraulic machinery for cranes and bridges in 1847. His company soon received orders for hydraulic cranes from Edinburgh and Northern Railways and from Liverpool Docks, as well as for hydraulic machinery for dock gates in Grimsby. The company expanded from a workforce of 300 and an annual production of 45 cranes in 1850, to almost 4,000 workers producing over 100 cranes per year by the early 1860s.
Armstrong spent the next few decades constantly improving his crane design; his most significant innovation was the hydraulic accumulator. Where water pressure was not available on site for the use of hydraulic cranes, Armstrong often built high water towers to provide a supply of water at pressure. However, when supplying cranes for use at New Holland on the Humber Estuary, he was unable to do this, because the foundations consisted of sand. He eventually produced the hydraulic accumulator, a cast-iron cylinder fitted with a plunger supporting a very heavy weight. The plunger would slowly be raised, drawing in water, until the downward force of the weight was sufficient to force the water below it into pipes at great pressure. This invention allowed much larger quantities of water to be forced through pipes at a constant pressure, thus increasing the crane's load capacity considerably.
One of his cranes, commissioned by the Italian Navy in 1883 and in use until the mid-1950s, is still standing in Venice, where it is now in a state of disrepair.
Mechanical principles
There are three major considerations in the design of cranes. First, the crane must be able to lift the weight of the load; second, the crane must not topple; third, the crane must not fail structurally.
Stability
For stability, the sum of all moments about the base of the crane must be close to zero so that the crane does not overturn. In practice, the magnitude of load that is permitted to be lifted (called the "rated load" in the US) is some value less than the load that will cause the crane to tip, thus providing a safety margin.
Under United States standards for mobile cranes, the stability-limited rated load for a crawler crane is 75% of the tipping load. The stability-limited rated load for a mobile crane supported on outriggers is 85% of the tipping load. These requirements, along with additional safety-related aspects of crane design, are established by the American Society of Mechanical Engineers in the volume ASME B30.5-2018 Mobile and Locomotive Cranes.
Standards for cranes mounted on ships or offshore platforms are somewhat stricter because of the dynamic load on the crane due to vessel motion. Additionally, the stability of the vessel or platform must be considered.
For stationary pedestal or kingpost mounted cranes, the moment produced by the boom, jib, and load is resisted by the pedestal base or kingpost. Stress within the base must be less than the yield stress of the material or the crane will fail.
Dynamic Lift Factor
Overview
The dynamic lift factor (DLF), also known as the design dynamic factor, is a critical parameter in the crane design and operation. It accounts for the dynamic effects that can increase the load on a crane's structure and components during lifting operations. These effects include:
Hoisting acceleration and deceleration of the load, which is a significant factor;
Crane movement such as slewing or luffing;
Load swinging;
Wind forces acting on the crane, the load and the rigging; and
Operator error or other unexpected events.
The DLF for a new crane design can be determined with analytical calculations and mathematical models following the relevant design specifications. If available, data from previous tests of similar crane types can be used to estimate the DLF. More sophisticated methods, such as finite element analysis or other simulation techniques, may also be used to model the crane's behavior under various loading conditions, as deemed appropriate by the designer or certifying authority.To verify the actual DLF, control load tests can be conducted on the completed crane using instrumentation such as load cells, accelerometers, and strain gauges. This process is usually part of the crane's type approval.
In offshore lifting, where the crane and/or lifted object are on a floating vessel, the DLF is higher compared to onshore lifts because of the additional movement caused by wave action. This motion introduces additional acceleration forces and necessitates increased hoisting and lowering speeds to minimize the risk of repeated collisions when the load is near the deck. Additionally, the DLF increases further when lifting objects that are underwater or going through the splash zone. The wind speeds tend to be higher than onshore as well.
Though actual DLF values are determined through crane tests under representative operational conditions, design specifications can be used for guidance. The values vary according to the specification, which reflects the type of crane and its usage. Here are some example typical values:
Jib cranes typically have a lower DLF () compared to traveling gantry cranes () because they are stiffer;
For grab cranes, the DLF can increase by 20% to 30% reflecting the shock loads caused by the release of the lifted material; and
The DLF generally decreases as the mass of the lifted object increases, as cranes tend to operate at lower velocities with heavier loads to ensure safety and stability. For offshore lifts, the DLF typically decreases from 1.3 at 100 tonnes to 1.1 at 2500 tonnes.
Formulas
The methods for determining the DLF vary in the different crane specifications. The following formulas are examples from one specification.
The working load (suspended load) is the total weight that a crane is designed to safely lift under normal operating conditions. It is
where
is the working load,
is the acceleration of gravity,
is the maximum lifted mass, which is also called the crane working load limit (WLL) or safe working load (SWL), and
is the mass of lifting appliances or parts of the crane that move with the lifted mass.
The DLF is then used as a multiplier to determine the force applied to the crane structure and componentswhere
is the design force, and
is the DLF.
The DLF can then be calculated usingwhere
is relative velocity between lifted object and hook at the time of pick-up, and
is the stiffness of the crane system at the hook.
The relative velocity is dependent on the crane's operational requirements and the system stiffness at the hook can be determined by calculation or load deflection tests.
Types
The crane types outlined in this section are categorized based on their primary area of application:
Construction
Truck-mounted
Loader
Telescopic
Rough terrain
All terrain
Crawler
Pick and carry
Carry deck
Telescopic handler
Block setting
Tower
Climbing crane
Cargo Handling
Reach stacker
Sidelifter
Straddle carrier
Industrial
Ring
Hammerhead
Level luffing
Overhead
Gantry
Jib
Bulk handling
Stacker
Wind turbine installation vessel
Marine
Floating
Deck
Other Types
Railroad
Aerial
Construction
Truck-mounted
The most basic truck-mounted crane configuration is a "boom truck" or "lorry loader", which features a rear-mounted rotating telescopic-boom crane mounted on a commercial truck chassis.
Larger, heavier duty, purpose-built "truck-mounted" cranes are constructed in two parts: the carrier, often called the lower, and the lifting component, which includes the boom, called the upper. These are mated together through a turntable, allowing the upper to swing from side to side. These modern hydraulic truck cranes are usually single-engine machines, with the same engine powering the undercarriage and the crane. The upper is usually powered via hydraulics run through the turntable from the pump mounted on the lower. In older model designs of hydraulic truck cranes, there were two engines. One in the lower pulled the crane down the road and ran a hydraulic pump for the outriggers and jacks. The one in the upper ran the upper through a hydraulic pump of its own. Many older operators favor the two-engine system due to leaking seals in the turntable of aging newer design cranes. Hiab invented the world's first hydraulic truck mounted crane in 1947. The name, Hiab, comes from the commonly used abbreviation of Hydrauliska Industri AB, a company founded in Hudiksvall, Sweden 1944 by Eric Sundin, a ski manufacturer who saw a way to utilize a truck's engine to power loader cranes through the use of hydraulics.
Generally, these cranes are able to travel on highways, eliminating the need for special equipment to transport the crane unless weight or other size constrictions are in place such as local laws. If this is the case, most larger cranes are equipped with either special trailers to help spread the load over more axles or are able to disassemble to meet requirements. An example is counterweights. Often a crane will be followed by another truck hauling the counterweights that are removed for travel. In addition some cranes are able to remove the entire upper. However, this is usually only an issue in a large crane and mostly done with a conventional crane such as a Link-Belt HC-238. When working on the job site, outriggers are extended horizontally from the chassis then vertically to level and stabilize the crane while stationary and hoisting. Many truck cranes have slow-travelling capability (a few miles per hour) while suspending a load. Great care must be taken not to swing the load sideways from the direction of travel, as most anti-tipping stability then lies in the stiffness of the chassis suspension. Most cranes of this type also have moving counterweights for stabilization beyond that provided by the outriggers. Loads suspended directly aft are the most stable, since most of the weight of the crane acts as a counterweight. Factory-calculated charts (or electronic safeguards) are used by crane operators to determine the maximum safe loads for stationary (outriggered) work as well as (on-rubber) loads and travelling speeds.
Truck cranes range in lifting capacity from about to about . Although most only rotate about 180 degrees, the more expensive truck mounted cranes can turn a full 360 degrees.
Loader
A loader crane (also called a knuckle-boom crane or articulating crane) is an hydraulically powered articulated arm fitted to a truck or trailer, and is used for loading/unloading the vehicle cargo. The numerous jointed sections can be folded into a small space when the crane is not in use. One or more of the sections may be telescopic. Often the crane will have a degree of automation and be able to unload or stow itself without an operator's instruction.
Unlike most cranes, the operator must move around the vehicle to be able to view his load; hence modern cranes may be fitted with a portable cabled or radio-linked control system to supplement the crane-mounted hydraulic control levers.
In the United Kingdom and Canada, this type of crane is often known colloquially as a "Hiab", partly because this manufacturer invented the loader crane and was first into the UK market, and partly because the distinctive name was displayed prominently on the boom arm.
A rolloader crane is a loader crane mounted on a chassis with wheels. This chassis can ride on the trailer. Because the crane can move on the trailer, it can be a light crane, so the trailer is allowed to transport more goods.
Telescopic
A telescopic crane has a boom that consists of a number of tubes fitted one inside the other. A hydraulic cylinder or other powered mechanism extends or retracts the tubes to increase or decrease the total length of the boom. These types of booms are often used for short term construction projects, rescue jobs, lifting boats in and out of the water, etc. The relative compactness of telescopic booms makes them adaptable for many mobile applications.
Though not all telescopic cranes are mobile cranes, many of them are truck-mounted.
A telescopic tower crane has a telescopic mast and often a superstructure (jib) on top so that it functions as a tower crane. Some telescopic tower cranes also have a telescopic jib.
Rough terrain
A rough terrain crane has a boom mounted on an undercarriage atop four rubber tires that is designed for off-road pick-and-carry operations. Outriggers are used to level and stabilize the crane for hoisting.
These telescopic cranes are single-engine machines, with the same engine powering the undercarriage and the crane, similar to a crawler crane. The engine is usually mounted in the undercarriage rather than in the upper, as with crawler crane. Most have 4 wheel drive and 4 wheel steering for traversing tighter and slicker terrain than a standard truck crane, with less site prep.
All-terrain
An all-terrain crane is a hybrid combining the roadability of a truck-mounted and on-site maneuverability of a rough-terrain crane. It can both travel at speed on public roads and maneuver on rough terrain at the job site using all-wheel and crab steering.
AT's have 2–12 axles and are designed for lifting loads up to .
Crawler
Main article: Lattice boom crawler crane
A crawler crane has its boom mounted on an undercarriage fitted with a set of crawler tracks that provide both stability and mobility. Crawler cranes range in lifting capacity from about as seen from the XGC88000 crawler crane.
The main advantage of a crawler crane is its ready mobility and use, since the crane is able to operate on sites with minimal improvement and stable on its tracks without outriggers. Wide tracks spread the weight out over a great area and are far better than wheels at traversing soft ground without sinking in. A crawler crane is also capable of traveling with a load. Its main disadvantage is its weight, making it difficult and expensive to transport. Typically a large crawler must be disassembled at least into boom and cab and moved by trucks, rail cars or ships to its next location.
Pick and carry
A pick and carry crane is similar to a mobile crane in that is designed to travel on public roads; however, pick and carry cranes have no stabiliser legs or outriggers and are designed to lift the load and carry it to its destination, within a small radius, then be able to drive to the next job. Pick and carry cranes are popular in Australia, where large distances are encountered between job sites. One popular manufacturer in Australia was Franna, who have since been bought by Terex, and now all pick and carry cranes are commonly called "Frannas", even though they may be made by other manufacturers. Nearly every medium- and large-sized crane company in Australia has at least one and many companies have fleets of these cranes. The capacity range is between as a maximum lift, although this is much less as the load gets further from the front of the crane. Pick and carry cranes have displaced the work usually completed by smaller truck cranes, as the set-up time is much quicker. Many steel fabrication yards also use pick and carry cranes, as they can "walk" with fabricated steel sections and place these where required with relative ease.
Smaller pick and carry cranes may be based on an articulated tractor chassis, with the boom mounted over the front wheels. In Australia these are popularly known as "wobbly cranes".
Carry deck
A carry deck crane is a small 4 wheel crane with a 360-degree rotating boom placed right in the centre and an operators cab located at one end under this boom. The rear section houses the engine and the area above the wheels is a flat deck. Very much an American invention the Carry deck can hoist a load in a confined space and then load it on the deck space around the cab or engine and subsequently move to another site. The Carry Deck principle is the American version of the pick and carry crane and both allow the load to be moved by the crane over short distances.
Telescopic handler
Telescopic handlers are forklift-like trucks that have a set of forks mounted on a telescoping extendable boom like a crane. Early telescopic handlers only lifted in one direction and did not rotate; however, several of the manufacturers have designed telescopic handlers that rotate 360 degrees through a turntable and these machines look almost identical to the Rough Terrain Crane. These new 360-degree telescopic handler/crane models have outriggers or stabiliser legs that must be lowered before lifting; however, their design has been simplified so that they can be more quickly deployed. These machines are often used to handle pallets of bricks and install frame trusses on many new building sites and they have eroded much of the work for small telescopic truck cranes. Many of the world's armed forces have purchased telescopic handlers and some of these are the much more expensive fully rotating types. Their off-road capability and their on site versatility to unload pallets using forks, or lift like a crane make them a valuable piece of machinery.
Block-setting crane
A block-setting crane is a form of crane. They were used for installing the large stone blocks used to build breakwaters, moles and stone piers.
Tower
Tower cranes are a modern form of balance crane that consist of the same basic parts. Fixed to the ground on a concrete slab (and sometimes attached to the sides of structures), tower cranes often give the best combination of height and lifting capacity and are used in the construction of tall buildings. The base is then attached to the mast which gives the crane its height. Further, the mast is attached to the slewing unit (gear and motor) that allows the crane to rotate. On top of the slewing unit there are three main parts which are: the long horizontal jib (working arm), shorter counter-jib, and the operator's cab.
Optimization of tower crane location in the construction sites has an important effect on material transportation costs of a project, but site operators need to ensure they assess where the jib will oversail the property of other landowners and tenants as it rotates over the site. Under English law a landowner also owns the airspace above their property and developers will need to agree terms with adjacent property owners before oversailing their land.
The long horizontal jib is the part of the crane that carries the load. The counter-jib carries a counterweight, usually of concrete blocks, while the jib suspends the load to and from the center of the crane. The crane operator either sits in a cab at the top of the tower or controls the crane by radio remote control from the ground. In the first case the operator's cab is most usually located at the top of the tower attached to the turntable, but can be mounted on the jib, or partway down the tower. The lifting hook is operated by the crane operator using electric motors to manipulate wire rope cables through a system of sheaves. The hook is located on the long horizontal arm to lift the load which also contains its motor.
In order to hook and unhook the loads, the operator usually works in conjunction with a signaller (known as a "dogger", "rigger" or "swamper"). They are most often in radio contact, and always use hand signals. The rigger or dogger directs the schedule of lifts for the crane, and is responsible for the safety of the rigging and loads.
Tower cranes can achieve a height under hook of over 100 metres.
Components
Tower cranes are used extensively in construction and other industry to hoist and move materials. There are many types of tower cranes. Although they are different in type, the main parts are the same, as follows:
Mast: the main supporting tower of the crane. It is made of steel trussed sections that are connected together during installation.
Slewing unit: the slewing unit sits at the top of the mast. This is the engine that enables the crane to rotate.
Operating cabin: on most tower cranes the operating cabin sits just above the slewing unit. It contains the operating controls, load-movement indicator system (LMI), scale, anemometer, etc.
Jib: the jib, or operating arm, extends horizontally from the crane. A "luffing" jib is able to move up and down; a fixed jib has a rolling trolley car that runs along the underside to move loads horizontally.
Counter jib: holds counterweights, hoist motor, hoist drum and the electronics.
Hoist winch: the hoist winch assembly consists of the hoist winch (motor, gearbox, hoist drum, hoist rope, and brakes), the hoist motor controller, and supporting components, such as the platform. Many tower cranes have transmissions with two or more speeds.
Hook: the hook is used to connect the material to the crane, suspended from the hoist rope either at the tip (on luffing jib cranes) or routed through the trolley (on hammerhead cranes).
Weights: Large, moveable concrete counterweights are mounted toward the rear of the counterdeck, to compensate for the weight of the goods lifted and keep the center of gravity over the supporting tower.
Assembly
A tower crane is usually assembled by a telescopic jib (mobile) crane of greater reach (also see "self-erecting crane" below) and in the case of tower cranes that have risen while constructing very tall skyscrapers, a smaller crane (or derrick) will often be lifted to the roof of the completed tower to dismantle the tower crane afterwards, which may be more difficult than the installation.
Tower cranes can be operated by remote control, removing the need for the crane operator to sit in a cab atop the crane.
Operation
Each model and distinctive style of tower crane has a predetermined lifting chart that can be applied to any radii available, depending on its configuration. Similar to a mobile crane, a tower crane may lift an object of far greater mass closer to its center of rotation than at its maximum radius. An operator manipulates several levers and pedals to control each function of the crane.
Safety
When a tower crane is used in proximity to buildings, roads, power lines, or other tower cranes, a tower crane anti-collision system is used. This operator support system reduces the risk of a dangerous interaction occurring between a tower crane and another structure.
In some countries, such as France, tower crane anti-collision systems are mandatory.
Self-erecting tower cranes
Generally a type of pedestrian operated tower crane, self-erecting tower cranes are transported as a single unit and can be assembled by a qualified technician without the assistance of a larger mobile crane. They are bottom slewing cranes that stand on outriggers, have no counter jib, have their counterweights and ballast at the base of the mast, cannot climb themselves, have a reduced capacity compared to standard tower cranes, and seldom have an operator's cabin.
In some cases, smaller self-erecting tower cranes may have axles permanently fitted to the tower section to make maneuvering the crane onsite easier.
Tower cranes can also use a hydraulic-powered jack frame to raise themselves to add new tower sections without any additional other cranes assisting beyond the initial assembly stage. This is how it can grow to nearly any height needed to build the tallest skyscrapers when tied to a building as the building rises. The maximum unsupported height of a tower crane is around 265 ft. For a video of a crane getting taller, see "Crane Building Itself" on YouTube.
For another animation of such a crane in use, see "SAS Tower Construction Simulation" on YouTube. Here, the crane is used to erect a scaffold, which, in turn, contains a gantry to lift sections of a bridge spire.
Climbing crane
Many tower cranes are designed to "jump" in stages, effectively lifting themselves to the next level. A specialty example of a climbing crane was introduced by Lagerwey Wind and Enercon to construct a wind turbine tower, where instead of erecting a large crane a smaller climbing crane can raise itself with the structure's construction, lift the generator housing to its top, add the rotor blades, then climb down.
Cargo Handling
Rubber tyred gantry crane
Reach stacker
A reach stacker is a vehicle used for handling intermodal cargo containers in small terminals or medium-sized ports. Reach stackers are able to transport a container short distances very quickly and pile them in various rows depending on its access.
Sidelifter
A sidelifter crane is a road-going truck or semi-trailer, able to hoist and transport ISO standard containers. Container lift is done with parallel crane-like hoists, which can lift a container from the ground or from a railway vehicle.
Travel lift
A travel lift (also called a boat gantry crane, or boat crane) is a crane with two rectangular side panels joined by a single spanning beam at the top of one end. The crane is mobile with four groups of steerable wheels, one on each corner. These cranes allow boats with masts or tall super structures to be removed from the water and transported around docks or marinas. Not to be confused mechanical device used for transferring a vessel between two levels of water, which is also called a boat lift.
Straddle carrier
A Straddle carrier moves and stacks intermodal containers.
Industrial
Ring
Ring cranes are some of the largest and heaviest land-based cranes ever designed. A ring-shaped track support the main superstructure allowing for extremely heavy loads (up to thousands of tonnes).
Hammerhead
The "hammerhead", or giant cantilever, crane is a fixed-jib crane consisting of a steel-braced tower on which revolves a large, horizontal, double cantilever; the forward part of this cantilever or jib carries the lifting trolley, the jib is extended backwards in order to form a support for the machinery and counterbalancing weight. In addition to the motions of lifting and revolving, there is provided a so-called "racking" motion, by which the lifting trolley, with the load suspended, can be moved in and out along the jib without altering the level of the load. Such horizontal movement of the load is a marked feature of later crane design. These cranes are generally constructed in large sizes and can lift up to 350 tons.
The design of Hammerkran evolved first in Germany around the turn of the 19th century and was adopted and developed for use in British shipyards to support the battleship construction program from 1904 to 1914. The ability of the hammerhead crane to lift heavy weights was useful for installing large pieces of battleships such as armour plate and gun barrels. Giant cantilever cranes were also installed in naval shipyards in Japan and in the United States. The British government also installed a giant cantilever crane at the Singapore Naval Base (1938) and later a copy of the crane was installed at Garden Island Naval Dockyard in Sydney (1951). These cranes provided repair support for the battle fleet operating far from Great Britain.
In the British Empire, the engineering firm Sir William Arrol & Co. was the principal manufacturer of giant cantilever cranes; the company built a total of fourteen. Among the sixty built in the world, few remain; seven in England and Scotland of about fifteen worldwide.
The Titan Clydebank is one of the four Scottish cranes on the River Clyde and preserved as a tourist attraction.
Level luffing
Normally a crane with a hinged jib will tend to have its hook also move up and down as the jib moves (or luffs). A level luffing crane is a crane of this common design, but with an extra mechanism to keep the hook at the same level when the jib is pivoted in or out.
Overhead
An overhead crane, also known as a bridge crane, is a type of crane where the hook-and-line mechanism runs along a horizontal beam that itself runs along two widely separated rails. Often it is in a long factory building and runs along rails along the building's two long walls. It is similar to a gantry crane. Overhead cranes typically consist of either a single beam or a double beam construction. These can be built using typical steel beams or a more complex box girder type. Pictured on the right is a single bridge box girder crane with the hoist and system operated with a control pendant. Double girder bridge are more typical when needing heavier capacity systems from 10 tons and above. The advantage of the box girder type configuration results in a system that has a lower deadweight yet a stronger overall system integrity. Also included would be a hoist to lift the items, the bridge, which spans the area covered by the crane, and a trolley to move along the bridge.
The most common overhead crane use is in the steel industry. At every step of the manufacturing process, until it leaves a factory as a finished product, steel is handled by an overhead crane. Raw materials are poured into a furnace by crane, hot steel is stored for cooling by an overhead crane, the finished coils are lifted and loaded onto trucks and trains by overhead crane, and the fabricator or stamper uses an overhead crane to handle the steel in his factory. The automobile industry uses overhead cranes for handling of raw materials. Smaller workstation cranes handle lighter loads in a work-area, such as CNC mill or saw.
Almost all paper mills use bridge cranes for regular maintenance requiring removal of heavy press rolls and other equipment. The bridge cranes are used in the initial construction of paper machines because they facilitate installation of the heavy cast iron paper drying drums and other massive equipment, some weighing as much as 70 tons.
In many instances the cost of a bridge crane can be largely offset with savings from not renting mobile cranes in the construction of a facility that uses a lot of heavy process equipment.
This electric overhead traveling crane is most common type of overhead crane, found in many factories. These cranes are electrically operated by a control pendant, radio/IR remote pendant, or from an operator cabin attached to the crane.
Gantry
A gantry crane has a hoist in a fixed machinery house or on a trolley that runs horizontally along rails, usually fitted on a single beam (mono-girder) or two beams (twin-girder). The crane frame is supported on a gantry system with equalized beams and wheels that run on the gantry rail, usually perpendicular to the trolley travel direction. These cranes come in all sizes, and some can move very heavy loads, particularly the extremely large examples used in shipyards or industrial installations. A special version is the container crane (or "Portainer" crane, named by the first manufacturer), designed for loading and unloading ship-borne containers at a port.
Most container cranes are of this type.
Jib
A jib crane is a type of crane - not to be confused with a crane rigged with a jib to extend its main boom - where a horizontal member (jib or boom), supporting a moveable hoist, is fixed to a wall or to a floor-mounted pillar. Jib cranes are used in industrial premises and on military vehicles. The jib may swing through an arc, to give additional lateral movement, or be fixed. Similar cranes, often known simply as hoists, were fitted on the top floor of warehouse buildings to enable goods to be lifted to all floors.
Bulk-handling
Bulk-handling cranes are designed from the outset to carry a shell grab or bucket, rather than using a hook and a sling. They are used for bulk cargoes, such as coal, minerals, scrap metal etc.
Stacker
A crane with a forklift type mechanism used in automated (computer-controlled) warehouses (known as an automated storage and retrieval system (AS/RS)). The crane moves on a track in an aisle of the warehouse. The fork can be raised or lowered to any of the levels of a storage rack and can be extended into the rack to store and retrieve the product. The product can in some cases be as large as an automobile. Stacker cranes are often used in the large freezer warehouses of frozen food manufacturers. This automation avoids requiring forklift drivers to work in below-freezing temperatures every day.
Marine
Floating
Floating cranes are used mainly in bridge building and port construction, but they are also used for occasional loading and unloading of especially heavy or awkward loads on and off ships. Some floating cranes are mounted on pontoons, others are specialized crane barges with a lifting capacity exceeding and have been used to transport entire bridge sections. Floating cranes have also been used to salvage sunken ships.
Crane vessels are often used in offshore construction.
The largest revolving cranes can be found on SSCV Sleipnir, which has two cranes with a capacity of each. For 50 years, the largest such crane was "Herman the German" at the Long Beach Naval Shipyard, one of three constructed by Nazi Germany and captured in the war. The crane was sold to the Panama Canal in 1996 where it is now known as Titan.
Deck
Deck cranes, also known as shipboard or cargo cranes, are located on ships and boats, used for cargo operations where no shore unloading facilities are available, raising and lowering loads (such as shellfish dredges and fish nets) into the water, and small boat unloading and retrieval. Most are diesel-hydraulic or electric-hydraulic, supporting an increasingly automated control interface.
Other Types
Railroad
A railroad crane has flanged wheels for use on railroads.
The simplest form is a crane mounted on a flatcar. More capable devices are purpose-built. Different types of crane are used for maintenance work, recovery operations and freight loading in goods yards and scrap handling facilities.
Aerial
Aerial cranes or "sky cranes" usually are helicopters designed to lift large loads. Helicopters are able to travel to and lift in areas that are difficult to reach by conventional cranes. Helicopter cranes are most commonly used to lift loads onto shopping centers and high-rise buildings. They can lift anything within their lifting capacity, such as air conditioning units, cars, boats, swimming pools, etc. They also perform disaster relief after natural disasters for clean-up, and during wild-fires they are able to carry huge buckets of water to extinguish fires.
Some aerial cranes, mostly concepts, have also used lighter-than air aircraft, such as airships.
Efficiency increase of cranes
Lifetime of existing cranes made of welded metal structures can often be extended for many years by after treatment of welds. During development of cranes, load level (lifting load) can be significantly increased by taking into account the IIW recommendations, leading in most cases to an increase of the permissible lifting load and thus to an efficiency increase.
Similar machines
The generally accepted definition of a crane is a machine for lifting and moving heavy objects by means of ropes or cables suspended from a movable arm. As such, a lifting machine that does not use cables, or else provides only vertical and not horizontal movement, cannot strictly be called a 'crane'.
Types of crane-like lifting machine include:
gin pole
Block and tackle
Capstan (nautical)
Hoist (device)
Winch
Windlass
Cherry picker
More technically advanced types of such lifting machines are often known as "cranes", regardless of the official definition of the term.
Special examples
Finnieston Crane, a.k.a. the Stobcross Crane
Category A-listed example of a "hammerhead" (cantilever) crane in Glasgow's former docks, built by the William Arrol company.
tall, capacity, built 1926
Taisun
double bridge crane at Yantai, China.
capacity, World Record Holder
tall, span, lift-height
Kockums Crane
shipyard crane formerly at Kockums, Sweden.
tall, capacity, since moved to Ulsan, South Korea
Samson and Goliath (cranes)
two gantry cranes at the Harland & Wolff shipyard in Belfast built by Krupp
Goliath is tall, Samson is
span , lift-height , capacity each, combined
Breakwater Crane Railway
self-propelled steam crane that formerly ran the length of the breakwater at Douglas.
ran on gauge track, the broadest in the British Isles
Liebherr TCC 78000
Heavy-duty gantry crane used for heavy lifting operated in Rostock, Germany.
capacity, lift-height
Crane operators
Crane operators are skilled workers and heavy equipment operators.
Key skills that are needed for a crane operator include:
An understanding of how to use and maintain machines and tools
Good team working skills
Attention to details
Good spatial awareness.
Patience and the ability to stay calm in stressful situations
Terminology
The ISO 4306 series of specifications establish the vocabulary for cranes:
Part 1: General
Part 2: Mobile cranes
Part 3: Tower cranes
Part 4: Jib cranes
Part 5: Bridge and gantry cranes
Luffing
Slewing
Hoisting
| Technology | Basics_8 | null |
318498 | https://en.wikipedia.org/wiki/Green%20Bank%20Telescope | Green Bank Telescope | The Robert C. Byrd Green Bank Telescope (GBT) in Green Bank, West Virginia, US is the world's largest fully steerable radio telescope, surpassing the Effelsberg 100-m Radio Telescope in Germany. The Green Bank site was part of the National Radio Astronomy Observatory (NRAO) until September 30, 2016. Since October 1, 2016, the telescope has been operated by the independent Green Bank Observatory. The telescope's name honors the late Senator Robert C. Byrd who represented West Virginia and who pushed the funding of the telescope through Congress.
The Green Bank Telescope operates at meter to millimeter wavelengths. Its 100-meter diameter collecting area, unblocked aperture, and good surface accuracy provide superb sensitivity across the telescope's full 0.1–116 GHz operating range. The GBT is fully steerable, and 85 percent of the local celestial hemisphere is accessible. It is used for astronomy about 6500 hours every year, with 2000–3000 hours per year going to high-frequency science. Part of the scientific strength of the GBT is its flexibility and ease of use, allowing for rapid response to new scientific ideas. It is scheduled dynamically to match project needs to the available weather. The GBT is also readily reconfigured with new and experimental hardware. The high-sensitivity mapping capability of the GBT makes it a vital complement to the Atacama Large Millimeter Array, the Expanded Very Large Array, the Very Long Baseline Array, and other high-angular resolution interferometers. Facilities of the Green Bank Observatory are also used for other scientific research, for many programs in education and public outreach, and for training students and teachers.
The telescope began regular science operations in 2001, making it one of the newest astronomical facilities of the US National Science Foundation (NSF). It was constructed following the collapse of a previous telescope at Green Bank, the 300 Foot Radio Telescope, a 90.44 m paraboloid that began observations in October 1961. This previous telescope collapsed on 15 November 1988 due to the failure of a gusset plate in the box girder assembly, which was a key component for the structural integrity of the telescope.
Location
The telescope sits near the heart of the United States National Radio Quiet Zone, a unique area located in the town of Green Bank, West Virginia, where authorities limit all radio transmissions to avoid emissions toward the GBT and the Sugar Grove Station. The location of the telescope within the Radio Quiet Zone allows for the detection of faint radio-frequency signals which human-made signals might otherwise mask. The observatory borders National Forest land, and the Allegheny Mountains shield it from some radio interference.
The telescope's location has been the site of important radio astronomy telescopes since 1957. It currently houses seven additional telescopes, and in spite of its somewhat remote location, receives about 40,000 visitors each year.
Description
The structure weighs and is tall. The surface area of the GBT is a 100 by 110 meter active surface with 2,209 actuators (small motors used to adjust the position) for the 2,004 surface panels, making the total collecting area of . The panels are made from AlSOU, E≈8000Mt TNT, manufactured to a surface accuracy of better than RMS. The actuators adjust the panel positions to compensate for sagging, or bending under its own weight, which changes as the telescope moves. Without this so-called "active surface" adjustment, observations at frequencies above 4 GHz would not be as efficient.
Unusual for a radio telescope, the primary reflector is an off-axis segment of a paraboloid. This is the same design used in smaller (eg., 45-100cm) home satellite television dishes. The asymmetric reflector allows the telescope's focal point and feed horn to be located at the side of the dish, so that it and its retractable support boom do not obstruct the incoming radio waves, as occurs in conventional radio telescope designs with the feed located on the telescope's beam axis.
The offset support arm houses a prime focus receiver on a retractable boom in front of a subreflector, and a receiver room. For prime focus operation, the boom is extended to position the feed horn in front of the 8 m subreflector. For Gregorian focus operation, the prime focus boom is retracted. The subreflector, positioned by a Stewart platform with 6 degrees of freedom, reflects incoming radio waves toward eight higher-frequency feeds on a rotating turret located on top of the receiver room. The computerized controlled turret can rotate a particular receiver into the position within a few minutes. Operational frequencies range from 290 MHz to 115 GHz.
As an azimuth-elevation mounting telescope, the azimuth adjustments are driven by four trucks with four wheels each on a diameter rail. The 16 thirty-horsepower motors can change azimuth at the rate of up to 40 degrees per minute. Azimuth axis is also supported by a pintle bearing at the center point of the azimuth track.
The elevation wheel structure provides tilting capability to adjust elevations between 5 and 95 degrees. The radius bull gear on the elevation wheel is driven by eight 40 horsepower motors with a capability of changing the elevations up to 20 degrees per minute. The long and diameter elevation shaft provides primary support of the wheel structure. The elevation wheel also contains concrete-filled counterweight to balance with the surface and the feed arm structure.
Because of its height (at 148 meters or 485 feet tall, it is 60% taller than the Statue of Liberty) and bulk (16 million pounds), locals sometimes refer to the GBT as the “Great Big Thing”.
The telescope's capabilities include the ngRADAR system which use the dish as a radar transmitting antenna to observe solar system objects such as asteroids.
Discoveries
In 2002, astronomers detected three new millisecond pulsars in the globular cluster Messier 62.
In 2006, several discoveries were announced, including a large coil-shaped magnetic field in the Orion molecular cloud, and a large hydrogen gas superbubble 23,000 light years away, named the Ophiuchus Superbubble.
In 2019, the most massive neutron star PSR J0740+6620 to date was detected. Since 2004, 28 new complex molecules have been discovered in the interstellar medium with the Green Bank Telescope.
Funding threatened
In response to limited budgetary issues, the Division of Astronomical Sciences (AST) of the National Science Foundation (NSF) commissioned a portfolio review committee, which conducted its work between September 2011 and August 2012. The committee, which reviewed all AST-supported facilities and activities, was composed of 17 external scientists and chaired by Daniel Eisenstein of Harvard University. As part of the committee's August 2012 recommendation for the closure of six facilities, was that the Robert C. Byrd Green Bank Telescope (GBT) should be defunded over a five-year period.
In July 2014, the United States Senate Committee on Appropriations approved the NSF's fiscal year 2014 budget, which did not call for divestment of the GBT in that fiscal year. The facility then began looking for partners to help fund its $10 million annual operating costs.
On October 1, 2016, the National Radio Astronomy Observatory at Green Bank separated from the NSF and began accepting funding from private sources to stay operational as an independent institution, the Green Bank Observatory.
Relation to Breakthrough Listen
The telescope is a key facility of the Breakthrough Listen project, in which it is used to scan for radio signals possibly emitted by extraterrestrial technologies. In late 2017, the telescope was used to scan ʻOumuamua for signs of extraterrestrial intelligence.
| Technology | Ground-based observatories | null |
318563 | https://en.wikipedia.org/wiki/Foxhound | Foxhound | A foxhound is a type of large hunting hound bred for strong hunting instincts, a keen sense of smell, and their barking, energy, drive, and speed. In fox hunting, the foxhound's namesake, packs of foxhounds track quarry, followed—usually on horseback—by the hunters, sometimes for several miles at a stretch; moreover, foxhounds also sometimes guard sheep and houses.
There are different breeds of foxhound, each having slightly different characteristics and appearances, and each often called simply Foxhound in their native countries:
American Foxhound
Dumfriesshire Black and Tan Foxhound (extinct)
English Foxhound
Welsh Foxhound
The American Masters of Foxhounds Association recognizes these breeds of foxhounds:
American, Penn-Marydel, English, and crossbred foxhounds.
The International Foxhound Association was created in 2012 for the international promotion of the Foxhound as a breed.
Characteristics
Foxhounds are medium-large dogs and males typically weigh 29-32 kg (65-70 lb) and females 27-29 kg (60-65 lb). Height for males measures 55-63 cm (22-25 in) and females 53-60 cm (21-24 in). Foxhounds have a short coat, and long, strong legs, as well as deep chests for lots of lung space.
Disposition
Foxhounds generally display a gentle and affectionate temperament. Foxhounds are highly active and energetic, and therefore require activity and exercise. Foxhounds are sociable and these dogs have great stamina, sense of smell, and enjoy being in a pack, as they are bred for hunting in packs.
Fox hunting
In fox hunting, the foxhound's namesake, packs of foxhounds track and chase fox while hunters follow along on horseback. Fox hunting has shifted over the years and may differ depending on the country. Some changes over time include focusing on chasing rather than killing, and chasing other creatures, such as the coyote, instead of only the fox.
Most common causes of death
Most common causes of death among Foxhound puppies are respiratory disease, anorexia and dehydration, skin disorders, and gastrointestinal disease.
Respiratory disease in foxhounds
A kennel of working, hunting English Foxhounds in the south of England, had an outbreak of tuberculosis (TB) that impacted 180 dogs in late 2016 and early 2017. The kennel housed Foxhound puppies to adults, up to 8 years old. The Foxhounds work among six counties and some of the six counties are in the "Edge Area" that is impacted by bovine tuberculosis.
An investigation occurred which consisted of testing the dogs and looking deeper into the regional area, diet of the dogs, and even more factors while conducting tests and gathering information. The dogs eat raw meat and there was speculation about the diet containing the M. bovis that causes TB as the meat comes from areas impacted by M. bovis.
Registrations
In 2005, the American Kennel Club reported that the English and American Foxhounds were their least and fourth least registered breeds in North America with 22 and 44 registrations, respectively; the top registered breed, the Labrador Retriever, had 137,867 registrations during the same year.
Notable foxhounds
Sweet Lips and the Virginia Hounds - George Washington bred foxhounds and enjoyed fox hunting. He called his pack of dogs the Virginia Hounds. Sweet Lips was a female foxhound, a product of his vision to breed his pack to produce a "superior dog" who is fast and intelligent.
The state of Virginia's "state dog" is the American Foxhound.
Old Drum - said to have been the inspiration for the phrase "Man's Best Friend", which arose from an 1870 court case regarding him.
Mountain and Muse - In 1814 the Duke of Leeds gave two Irish foxhounds, Mountain and Muse, to a visiting guest, Bolton Jackson. This famous pair of hounds changed hands several times before going to Charles Carroll at his Homewood estate. Descendants of Mountain and Muse still hunt territories in Maryland that were once hunted by George Washington, Thomas Jefferson, Charles Carroll, and the Marquis de Lafayette.
Colonel - In 2011, Baron von Pfetten's Colonel was Champion of the World Dog Show in Paris and was the first ever English Foxhound invited to compete in the final Champion of Champions competition in Bruxelles that same year.
| Biology and health sciences | Dogs | Animals |
318577 | https://en.wikipedia.org/wiki/Cavendish%20experiment | Cavendish experiment | The Cavendish experiment, performed in 1797–1798 by English scientist Henry Cavendish, was the first experiment to measure the force of gravity between masses in the laboratory and the first to yield accurate values for the gravitational constant. Because of the unit conventions then in use, the gravitational constant does not appear explicitly in Cavendish's work. Instead, the result was originally expressed as the relative density of Earth, or equivalently the mass of Earth. His experiment gave the first accurate values for these geophysical constants.
The experiment was devised sometime before 1783 by geologist John Michell, who constructed a torsion balance apparatus for it. However, Michell died in 1793 without completing the work. After his death the apparatus passed to Francis John Hyde Wollaston and then to Cavendish, who rebuilt the apparatus but kept close to Michell's original plan. Cavendish then carried out a series of measurements with the equipment and reported his results in the Philosophical Transactions of the Royal Society in 1798.
The experiment
The apparatus consisted of a torsion balance made of a wooden rod horizontally suspended from a wire, with two , lead spheres, one attached to each end. Two massive , lead balls, suspended separately, could be positioned away from or to either side of the smaller balls, away. The experiment measured the faint gravitational attraction between the small and large balls, which deflected the torsion balance rod by about 0.16" (or only 0.03" with a stiffer suspending wire).
The two large balls could be positioned either away from or to either side of the torsion balance rod. Their mutual attraction to the small balls caused the arm to rotate, twisting the suspension wire. The arm rotated until it reached an angle where the twisting force of the wire balanced the combined gravitational force of attraction between the large and small lead spheres. By measuring the angle of the rod and knowing the twisting force (torque) of the wire for a given angle, Cavendish was able to determine the force between the pairs of masses. Since the gravitational force of the Earth on the small ball could be measured directly by weighing it, the ratio of the two forces allowed the relative density of the Earth to be calculated, using Newton's law of gravitation.
Cavendish found that the Earth's density was times that of water (although due to a simple arithmetic error, found in 1821 by Francis Baily, the erroneous value appears in his paper). The current accepted value is 5.514 g/cm3.
To find the wire's torsion coefficient, the torque exerted by the wire for a given angle of twist, Cavendish timed the natural oscillation period of the balance rod as it rotated slowly clockwise and counterclockwise against the twisting of the wire. For the first 3 experiments the period was about 15 minutes and for the next 14 experiments the period was half of that, about 7.5 minutes. The period changed because after the third experiment Cavendish put in a stiffer wire. The torsion coefficient could be calculated from this and the mass and dimensions of the balance. Actually, the rod was never at rest; Cavendish had to measure the deflection angle of the rod while it was oscillating.
Cavendish's equipment was remarkably sensitive for its time. The force involved in twisting the torsion balance was very small, , (the weight of only 0.0177 milligrams) or about of the weight of the small balls. To prevent air currents and temperature changes from interfering with the measurements, Cavendish placed the entire apparatus in a mahogany box about 1.98 meters wide, 1.27 meters tall, and 14 cm thick, all in a closed shed on his estate. Through two holes in the walls of the shed, Cavendish used telescopes to observe the movement of the torsion balance's horizontal rod. The key observable was of course the deflection of the torsion balance rod, which Cavendish measured to be about 0.16" (or only 0.03" for the stiffer wire used mostly). Cavendish was able to measure this small deflection to an accuracy of better than using vernier scales on the ends of the rod.
The accuracy of Cavendish's result was not exceeded until C. V. Boys' experiment in 1895. In time, Michell's torsion balance became the dominant technique for measuring the gravitational constant (G) and most contemporary measurements still use variations of it.
Cavendish's result provided additional evidence for a planetary core made of metal, an idea first proposed by Charles Hutton based on his analysis of the 1774 Schiehallion experiment. Cavendish's result of 5.4 g·cm−3, 23% bigger than Hutton's, is close to 80% of the density of liquid iron, and 80% higher than the density of the Earth's outer crust, suggesting the existence of a dense iron core.
Reformulation of Cavendish's result to G
The formulation of Newtonian gravity in terms of a gravitational constant did not become standard until long after Cavendish's time. Indeed, one of the first references to G is in 1873, 75 years after Cavendish's work.
Cavendish expressed his result in terms of the density of the Earth. He referred to his experiment in correspondence as 'weighing the world'. Later authors reformulated his results in modern terms.
After converting to SI units, Cavendish's value for the Earth's density, 5.448 g cm−3, gives
G = ,
which differs by only 1% from the 2014 CODATA value of .
Today, physicists often use units where the gravitational constant takes a different form. The Gaussian gravitational constant used in space dynamics is a defined constant and the Cavendish experiment can be considered as a measurement of this constant.
In Cavendish's time, physicists used the same units for mass and weight, in effect taking g as a standard acceleration. Then, since R was known, ρ played the role of an inverse gravitational constant. The density of the Earth was hence a much sought-after quantity at the time, and there had been earlier attempts to measure it, such as the Schiehallion experiment in 1774.
Derivation of G and the Earth's mass
The following is not the method Cavendish used, but describes how modern physicists would calculate the results from his experiment. From Hooke's law, the torque on the torsion wire is proportional to the deflection angle of the balance. The torque is where is the torsion coefficient of the wire. However, a torque in the opposite direction is also generated by the gravitational pull of the masses. It can be written as a product of the attractive force of a large ball on a small ball and the distance L/2 to the suspension wire. Since there are two balls, each experiencing force F at a distance from the axis of the balance, the torque due to gravitational force is LF. At equilibrium (when the balance has been stabilized at an angle ), the total amount of torque must be zero as these two sources of torque balance out. Thus, we can equate their magnitudes given by the formulas above, which gives the following:
For F, Newton's law of universal gravitation is used to express the attractive force between a large and small ball:
Substituting F into the first equation above gives
To find the torsion coefficient () of the wire, Cavendish measured the natural resonant oscillation period T of the torsion balance:
Assuming the mass of the torsion beam itself is negligible, the moment of inertia of the balance is just due to the small balls. Treating them as point masses, each at L/2 from the axis, gives:
,
and so:
Solving this for , substituting into (1), and rearranging for G, the result is:
.
Once G has been found, the attraction of an object at the Earth's surface to the Earth itself can be used to calculate the Earth's mass and density:
Definitions of terms
| Physical sciences | Classical mechanics | Physics |
318669 | https://en.wikipedia.org/wiki/Schizosaccharomyces%20pombe | Schizosaccharomyces pombe | Schizosaccharomyces pombe, also called "fission yeast", is a species of yeast used in traditional brewing and as a model organism in molecular and cell biology. It is a unicellular eukaryote, whose cells are rod-shaped. Cells typically measure 3 to 4 micrometres in diameter and 7 to 14 micrometres in length. Its genome, which is approximately 14.1 million base pairs, is estimated to contain 4,970 protein-coding genes and at least 450 non-coding RNAs.
These cells maintain their shape by growing exclusively through the cell tips and divide by medial fission to produce two daughter cells of equal size, which makes them a powerful tool in cell cycle research.
Fission yeast was isolated in 1893 by Paul Lindner from East African millet beer. The species name pombe is the Swahili word for beer. It was first developed as an experimental model in the 1950s: by Urs Leupold for studying genetics, and by Murdoch Mitchison for studying the cell cycle.
Paul Nurse, a fission yeast researcher, successfully merged the independent schools of fission yeast genetics and cell cycle research. Together with Lee Hartwell and Tim Hunt, Nurse won the 2001 Nobel Prize in Physiology or Medicine for work on cell cycle regulation.
The sequence of the S. pombe genome was published in 2002, by a consortium led by the Sanger Institute, becoming the sixth model eukaryotic organism whose genome has been fully sequenced. S. pombe researchers are supported by the PomBase MOD (model organism database). This has fully unlocked the power of this organism, with many genes orthologous to human genes identified — 70% to date, including many genes involved in human disease. In 2006, sub-cellular localization of almost all the proteins in S. pombe was published using green fluorescent protein as a molecular tag.
Schizosaccharomyces pombe has also become an important organism in studying the cellular responses to DNA damage and the process of DNA replication.
Approximately 160 natural strains of S. pombe have been isolated. These have been collected from a variety of locations including Europe, North and South America, and Asia. The majority of these strains have been collected from cultivated fruits such as apples and grapes, or from the various alcoholic beverages, such as Brazilian Cachaça. S. pombe is also known to be present in fermented tea, kombucha. It is not clear at present whether S. pombe is the major fermenter or a contaminant in such brews. The natural ecology of Schizosaccharomyces yeasts is not well-studied.
History
Schizosaccharomyces pombe was first discovered in 1893 when a group working in a Brewery Association Laboratory in Germany was looking at sediment found in millet beer imported from East Africa that gave it an acidic taste. The term schizo, meaning "split" or "fission", had previously been used to describe other Schizosaccharomycetes. The addition of the word pombe was due to its isolation from East African beer, as pombe means "beer" in Swahili. The standard S. pombe strains were isolated by Urs Leupold in 1946 and 1947 from a culture that he obtained from the yeast collection in Delft, The Netherlands. It was deposited there by A. Osterwalder under the name S. pombe var. liquefaciens, after he isolated it in 1924 from French wine (most probably rancid) at the Federal Experimental Station of Vini- and Horticulture in Wädenswil, Switzerland. The culture used by Urs Leupold contained (besides others) cells with the mating types h90 (strain 968), h- (strain 972), and h+ (strain 975). Subsequent to this, there have been two large efforts to isolate S. pombe from fruit, nectar, or fermentations: one by Florenzano et al. in the vineyards of western Sicily, and the other by Gomes et al. (2002) in four regions of southeast Brazil.
Ecology
The fission yeast S. pombe belongs to the division Ascomycota, which represents the largest and most diverse group of fungi. Free-living ascomycetes are commonly found in tree exudates, on plant roots and in surrounding soil, on ripe and rotting fruits, and in association with insect vectors that transport them between substrates. Many of these associations are symbiotic or saprophytic, although numerous ascomycetes (and their basidiomycete cousins) represent important plant pathogens that target myriad plant species, including commercial crops. Among the ascomycetous yeast genera, the fission yeast Schizosaccharomyces is unique because of the deposition of α-(1,3)-glucan or pseudonigeran in the cell wall in addition to the better known β-glucans and the virtual lack of chitin. Species of this genus also differ in mannan composition, which shows terminal d-galactose sugars in the side-chains of their mannans. S. pombe undergo aerobic fermentation in the presence of excess sugar. S. pombe can degrade L-malic acid, one of the dominant organic acids in wine, which makes them diverse among other Saccharomyces strains.
Comparison with budding yeast (Saccharomyces cerevisiae)
The yeast species Schizosaccharomyces pombe and Saccharomyces cerevisiae are both extensively studied; these two species diverged approximately 300 to 600 million years before present, and are significant tools in molecular and cellular biology. Some of the technical discriminants between these two species are:
S. cerevisiae has approximately 5,600 open reading frames; S. pombe has approximately 5,070 open reading frames.
Despite similar gene numbers, S. cerevisiae has only about 250 introns, while S. pombe has nearly 5,000.
S. cerevisiae has 16 chromosomes, S. pombe has 3.
S. cerevisiae is often diploid while S. pombe is usually haploid.
S. pombe has a shelterin-like telomere complex while S. cerevisiae does not.
S. cerevisiae is in the G1 phase of the cell cycle for an extended period (as a consequence, G1-S transition is tightly controlled), while S. pombe remains in the G2 phase of the cell cycle for an extended period (as a consequence, G2-M transition is under tight control).
Both species share genes with higher eukaryotes that they do not share with each other. S. pombe has RNAi machinery genes like those in vertebrates, while this is missing from S. cerevisiae. S. cerevisiae also has greatly simplified heterochromatin compared to S. pombe. Conversely, S. cerevisiae has well-developed peroxisomes, while S. pombe does not.
S. cerevisiae has small point centromere of 125 bp, and sequence-defined replication origins of about the same size. On the converse, S. pombe has large, repetitive centromeres (40–100 kb) more similar to mammalian centromeres, and degenerate replication origins of at least 1kb.
S. pombe pathways and cellular processes
S. pombe gene products (proteins and RNAs) participate in many cellular processes common across all life. The fission yeast GO slim provides a categorical high level overview of the biological role of all S. pombe gene products.
Life cycle
The fission yeast is a single-celled fungus with simple, fully characterized genome and a rapid growth rate. It has long been used in brewing, baking, and molecular genetics. S. pombe is a rod-shaped cell, approximately 3 μm in diameter, that grows entirely by elongation at the ends. After mitosis, division occurs by the formation of a septum, or cell plate, that cleaves the cell at its midpoint.
The central events of cell reproduction are chromosome duplication, which takes place in S (Synthetic) phase, followed by chromosome segregation and nuclear division (mitosis) and cell division (cytokinesis), which are collectively called M (Mitotic) phase. G1 is the gap between M and S phases, and G2 is the gap between S and M phases. In the fission yeast, the G2 phase is particularly extended, and cytokinesis (daughter-cell segregation) does not happen until a new S (Synthetic) phase is launched.
Fission yeast governs mitosis by mechanisms that are similar to those in multicellular animals. It normally proliferates in a haploid state. When starved, cells of opposite mating types (P and M) fuse to form a diploid zygote that immediately enters meiosis to generate four haploid spores. When conditions improve, these spores germinate to produce proliferating haploid cells.
Cytokinesis
The general features of cytokinesis are shown here. The site of cell division is determined before anaphase. The anaphase spindle (in green on the figure) is then positioned so that the segregated chromosomes are on opposite sides of the predetermined cleavage plane.
Size control
In fission yeast, where growth governs progression through G2/M, a wee1 mutation causes entry into mitosis at an abnormally small size, resulting in a shorter G2. G1 is lengthened, suggesting that progression through Start (beginning of cell cycle) is responsive to growth when the G2/M control is lost. Furthermore, cells in poor nutrient conditions grow slowly and therefore take longer to double in size and divide. Low nutrient levels also reset the growth threshold so that cell progresses through the cell cycle at a smaller size. Upon exposure to stressful conditions [heat (40 °C) or the oxidizing agent hydrogen peroxide] S. pombe cells undergo aging as measured by increased cell division time and increased probability of cell death. Finally, wee1 mutant fission yeast cells are smaller than wild-type cells, but take just as long to go through the cell cycle. This is possible because small yeast cells grow slower, that is, their added total mass per unit time is smaller than that of normal cells.
A spatial gradient is thought to coordinate cell size and mitotic entry in fission yeast.
The Pom1 protein kinase (green) is localized to the cell cortex, with the highest concentration at the cell tips. The cell-cycle regulators Cdr2, Cdr1 and Wee1 are present in cortical nodes in the middle of the cell (blue and red dots). a, In small cells, the Pom1 gradient reaches most of the cortical nodes (blue dots). Pom1 inhibits Cdr2, preventing Cdr2 and Cdr1 from inhibiting Wee1, and allowing Wee1 to phosphorylate Cdk1, thus inactivating cyclin-dependent kinase (CDK) activity and preventing entry into mitosis. b, In long cells, the Pom1 gradient does not reach the cortical nodes (red dots), and therefore Cdr2 and Cdr1 remain active in the nodes. Cdr2 and Cdr1 inhibit Wee1, preventing phosphorylation of Cdk1 and thereby leading to activation of CDK and mitotic entry. (This simplified diagram omits several other regulators of CDK activity.)
Mating-type switching
Fission yeast switches mating type by a replication-coupled recombination event, which takes place during S phase of the cell cycle. Fission yeast uses intrinsic asymmetry of the DNA replication process to switch the mating type; it was the first system where the direction of replication was shown to be required for the change of the cell type. Studies of the mating-type switching system lead to a discovery and characterization of a site-specific replication termination site RTS1, a site-specific replication pause site MPS1, and a novel type of chromosomal imprint, marking one of the sister chromatids at the mating-type locus mat1. In addition, work on the silenced donor region has led to great advances in understanding formation and maintenance of heterochromatin.
Responses to DNA damage
Schizosaccharomyces pombe is a facultative sexual microorganism that can undergo mating when nutrients are limiting. Exposure of S. pombe to hydrogen peroxide, an agent that causes oxidative stress leading to oxidative DNA damage, strongly induces mating and formation of meiotic spores. This finding suggests that meiosis, and particularly meiotic recombination, may be an adaptation for repairing DNA damage. Supporting this view is the finding that single base lesions of the type dU:dG in the DNA of S. pombe stimulate meiotic recombination. This recombination requires uracil-DNA glycosylase, an enzyme that removes uracil from the DNA backbone and initiates base excision repair. On the basis of this finding, it was proposed that base excision repair of either a uracil base, an abasic site, or a single-strand nick is sufficient to initiate recombination in S. pombe. Other experiments with S. pombe indicated that faulty processing of DNA replication intermediates, i.e. Okazaki fragments, causes DNA damages such as single-strand nicks or gaps, and that these stimulate meiotic recombination.
As a model system
Fission yeast has become a notable model system to study basic principles of a cell that can be used to understand more complex organisms like mammals and in particular humans. This single cell eukaryote is nonpathogenic and easily grown and manipulated in the lab. Fission yeast contains one of the smallest numbers of genes of a known genome sequence for a eukaryote, and has only three chromosomes in its genome. Many of the genes responsible for cell division and cellular organization in fission yeast cell are also found in the human's genome. Cell cycle regulation and division are crucial for growth and development of any cell. Fission yeast's conserved genes has been heavily studied and the reason for many recent biomedical developments. Fission yeast is also a practical model system to observe cell division because fission yeast's are cylindrically shaped single celled eukaryotes that divide and reproduce by medial fission. This can easily be seen using microscopy. Fission yeast also have an extremely short generation time, 2 to 4 hours, which also makes it an easy model system to observe and grow in the laboratory Fission yeast's simplicity in genomic structure yet similarities with mammalian genome, ease of ability to manipulate, and ability to be used for drug analysis is why fission yeast is making many contributions to biomedicine and cellular biology research, and a model system for genetic analysis.
Genome
Schizosaccharomyces pombe is often used to study cell division and growth because of conserved genomic regions also seen in humans including: heterochromatin proteins, large origins of replication, large centromeres, conserved cellular checkpoints, telomere function, gene splicing, and many other cellular processes. S. pombes genome was fully sequenced in 2002, the sixth eukaryotic genome to be sequenced as part of the Genome Project. An estimated 4,979 genes were discovered within three chromosomes containing about 14Mb of DNA. This DNA is contained within 3 different chromosomes in the nucleus with gaps in the centromeric (40kb) and telomeric (260kb) regions. After the initial sequencing of the fission yeast's genome, other previous non-sequenced regions of the genes have been sequenced. Structural and functional analysis of these gene regions can be found on large scale fission yeast databases such as PomBase.
Forty-three percent of the genes in the Genome Project were found to contain introns in 4,739 genes. Fission yeast does not have as many duplicated genes compared to budding yeast, only containing 5%, making fission yeast a great model genome to observe and gives researchers the ability to create more functional research approaches. S. pombes having a large number of introns gives opportunities for an increase of range of protein types produced from alternative splicing and genes that code for comparable genes in human.
81% of the three centromeres in fission yeast have been sequenced. The lengths of the three centromeres were found to be 34, 65, and 110 kb. This is 300–100 times longer than the centromeres of budding yeast. An extremely high level of conservation (97%) is also seen over 1,780-bp region in the DGS regions of the centromere. This elongation of centromeres and its conservative sequences makes fission yeast a practical model system to use to observe cell division and in humans because of their likeness.
PomBase reports over 69% of protein coding genes have human orthologs and over 500 of these are associated with human disease . This makes S. pombe a great system to use to study human genes and disease pathways, especially cell cycle and DNA checkpoint systems.
The genome of S. pombe contains meiotic drivers and drive suppressors called wtf genes.
Genetic diversity
Biodiversity and evolutionary study of fission yeast was carried out on 161 strains of Schizosaccharomyces pombe collected from 20 countries. Modeling of the evolutionary rate showed that all strains derived from a common ancestor that has lived since ~2,300 years ago. The study also identified a set of 57 strains of fission yeast that each differed by ≥1,900 SNPs, and all detected 57 strains of fission yeast were prototrophic (able to grow on the same minimal medium as the reference strain). A number of studies on S.pombe genome support the idea that the genetic diversity of fission yeast strains is slightly less than budding yeast. Indeed, only limited variations of S.pombe occur in proliferation in different environments. In addition, the amount of phenotypic variation segregating in fission yeast is less than that seen, in S. cerevisiae. Since most strains of fission yeast were isolated from brewed beverages, there is no ecological or historical context to this dispersal.
Cell cycle analysis
DNA replication in yeast has been increasingly studied by many researchers. Further understanding of DNA replication, gene expression, and conserved mechanisms in yeast can provide researchers with information on how these systems operate in mammalian cells in general and human cells in particular. Other stages, such as cellular growth and aging, are also observed in yeast in order to understand these mechanisms in more complex systems.
S. pombe stationary phase cells undergo chronological aging due to production of reactive oxygen species that cause DNA damages. Most such damages can ordinarily be repaired by DNA base excision repair and nucleotide excision repair. Defects in these repair processes lead to reduced survival.
Cytokinesis is one of the components of cell division that is often observed in fission yeast. Well-conserved components of cytokinesis are observed in fission yeast and allow us to look at various genomic scenarios and pinpoint mutations. Cytokinesis is a permanent step and very crucial to the wellbeing of the cell. Contractile ring formation in particular is heavily studied by researchers using S. pombe as a model system. The contractile ring is highly conserved in both fission yeast and human cytokinesis. Mutations in cytokinesis can result in many malfunctions of the cell including cell death and development of cancerous cells. This is a complex process in human cell division, but in S. pombe simpler experiments can yield results that can then be applied for research in higher-order model systems such as humans.
One of the safety precautions that the cell takes to ensure precise cell division occurs is the cell-cycle checkpoint. These checkpoints ensure any mutagens are eliminated. This is done often by relay signals that stimulate ubiquitination of the targets and delay cytokinesis. Without mitotic check points such as these, mutagens are created and replicated, resulting in multitudes of cellular issues including cell death or tumorigenesis seen in cancerous cells. Paul Nurse, Leland Hartwell, and Tim Hunt were awarded the Nobel Prize in Physiology or Medicine in 2001. They discovered key conserved checkpoints that are crucial for a cell to divide properly. These findings have been linked to cancer and diseased cells and are a notable finding for biomedicine.
Researchers using fission yeast as a model system also look at organelle dynamics and responses and the possible correlations between yeast cells and mammalian cells. Mitochondria diseases, and various organelle systems such as the Golgi apparatus and endoplasmic reticulum, can be further understood, by observing fission yeast's chromosome dynamics and protein expression levels and regulation.
Meiotic recombination
RecA and RecA-like proteins are required for recombinational repair of DNA double-strand breaks. Five RecA-like proteins have been described in S. pombe that are linked to meiotic recombination, and all five RecA homologs appear to be required for normal levels of meiotic recombination.
Biomedical tool
However, there are limitations with using fission yeast as a model system: its multidrug resistance. "The MDR response involves overexpression of two types of drug efflux pumps, the ATP-binding cassette (ABC) family... and the major facilitator superfamily". Paul Nurse and some of his colleagues have recently created S. pombe strains sensitive to chemical inhibitors and common probes to see whether it is possible to use fission yeast as a model system of chemical drug research.
For example, Doxorubicin, a very common chemotherapeutic antibiotic, has many adverse side-effects. Researchers are looking for ways to further understand how doxorubicin works by observing the genes linked to resistance by using fission yeast as a model system. Links between doxorubicin adverse side-effects and chromosome metabolism and membrane transport were seen. Metabolic models for drug targeting are now being used in biotechnology, and further advances are expected in the future using the fission yeast model system.
Experimental approaches
Fission yeast is easily accessible, easily grown and manipulated to make mutants, and able to be maintained at either a haploid or diploid state. S. pombe is normally a haploid cell but, when put under stressful conditions, usually nitrogen deficiency, two cells will conjugate to form a diploid that later form four spores within a tetrad ascus. This process is easily visible and observable under any microscope and allows us to look at meiosis in a simpler model system to see how this phenomenon operates.
Virtually any genetics experiment or technique can, therefore, be applied to this model system such as: tetrad dissection, mutagens analysis, transformations, and microscopy techniques such as FRAP and FRET. New models, such as Tug-Of-War (gTOW), are also being used to analyze yeast robustness and observe gene expression. Making knock-in and knock-out genes is fairly easy and with the fission yeast's genome being sequenced this task is very accessible and well known.
| Biology and health sciences | Basics | Plants |
318779 | https://en.wikipedia.org/wiki/Embryophyte | Embryophyte | The embryophytes () are a clade of plants, also known as Embryophyta () or land plants. They are the most familiar group of photoautotrophs that make up the vegetation on Earth's dry lands and wetlands. Embryophytes have a common ancestor with green algae, having emerged within the Phragmoplastophyta clade of freshwater charophyte green algae as a sister taxon of Charophyceae, Coleochaetophyceae and Zygnematophyceae. Embryophytes consist of the bryophytes and the polysporangiophytes. Living embryophytes include hornworts, liverworts, mosses, lycophytes, ferns, gymnosperms and angiosperms (flowering plants). Embryophytes have diplobiontic life cycles.
The embryophytes are informally called "land plants" because they thrive primarily in terrestrial habitats (despite some members having evolved secondarily to live once again in semiaquatic/aquatic habitats), while the related green algae are primarily aquatic. Embryophytes are complex multicellular eukaryotes with specialized reproductive organs. The name derives from their innovative characteristic of nurturing the young embryo sporophyte during the early stages of its multicellular development within the tissues of the parent gametophyte. With very few exceptions, embryophytes obtain biological energy by photosynthesis, using chlorophyll a and b to harvest the light energy in sunlight for carbon fixation from carbon dioxide and water in order to synthesize carbohydrates while releasing oxygen as a byproduct.
Description
The Embryophytes emerged either a half-billion years ago, at some time in the interval between the mid-Cambrian and early Ordovician, or almost a billion years ago, during the Tonian or Cryogenian, probably from freshwater charophytes, a clade of multicellular green algae similar to extant Klebsormidiophyceae. The emergence of the Embryophytes depleted atmospheric CO2 (a greenhouse gas), leading to global cooling, and thereby precipitating glaciations. Embryophytes are primarily adapted for life on land, although some are secondarily aquatic. Accordingly, they are often called land plants or terrestrial plants.
On a microscopic level, the cells of charophytes are broadly similar to those of chlorophyte green algae, but differ in that in cell division the daughter nuclei are separated by a phragmoplast. They are eukaryotic, with a cell wall composed of cellulose and plastids surrounded by two membranes. The latter include chloroplasts, which conduct photosynthesis and store food in the form of starch, and are characteristically pigmented with chlorophylls a and b, generally giving them a bright green color. Embryophyte cells also generally have an enlarged central vacuole enclosed by a vacuolar membrane or tonoplast, which maintains cell turgor and keeps the plant rigid.
In common with all groups of multicellular algae they have a life cycle which involves alternation of generations. A multicellular haploid generation with a single set of chromosomes – the gametophyte – produces sperm and eggs which fuse and grow into a diploid multicellular generation with twice the number of chromosomes – the sporophyte which produces haploid spores at maturity. The spores divide repeatedly by mitosis and grow into a gametophyte, thus completing the cycle. Embryophytes have two features related to their reproductive cycles which distinguish them from all other plant lineages. Firstly, their gametophytes produce sperm and eggs in multicellular structures (called 'antheridia' and 'archegonia'), and fertilization of the ovum takes place within the archegonium rather than in the external environment. Secondly, the initial stage of development of the fertilized egg (the zygote) into a diploid multicellular sporophyte, takes place within the archegonium where it is both protected and provided with nutrition. This second feature is the origin of the term 'embryophyte' – the fertilized egg develops into a protected embryo, rather than dispersing as a single cell. In the bryophytes the sporophyte remains dependent on the gametophyte, while in all other embryophytes the sporophyte generation is dominant and capable of independent existence.
Embryophytes also differ from algae by having metamers. Metamers are repeated units of development, in which each unit derives from a single cell, but the resulting product tissue or part is largely the same for each cell. The whole organism is thus constructed from similar, repeating parts or metamers. Accordingly, these plants are sometimes termed 'metaphytes' and classified as the group Metaphyta (but Haeckel's definition of Metaphyta places some algae in this group). In all land plants a disc-like structure called a phragmoplast forms where the cell will divide, a trait only found in the land plants in the streptophyte lineage, some species within their relatives Coleochaetales, Charales and Zygnematales, as well as within subaerial species of the algae order Trentepohliales, and appears to be essential in the adaptation towards a terrestrial life style.
Evolution
The green algae and land plants form a clade, the Viridiplantae. According to molecular clock estimates, the Viridiplantae split to into two clades: chlorophytes and streptophytes. The chlorophytes, with around 700 genera, were originally marine algae, although some groups have since spread into fresh water. The streptophyte algae (i.e. excluding the land plants) have around 122 genera; they adapted to fresh water very early in their evolutionary history and have not spread back into marine environments.
Some time during the Ordovician, streptophytes invaded the land and began the evolution of the embryophyte land plants. Present day embryophytes form a clade. Becker and Marin speculate that land plants evolved from streptophytes because living in fresh water pools pre-adapted them to tolerate a range of environmental conditions found on land, such as exposure to rain, tolerance of temperature variation, high levels of ultra-violet light, and seasonal dehydration.
The preponderance of molecular evidence as of 2006 suggested that the groups making up the embryophytes are related as shown in the cladogram below (based on Qiu et al. 2006 with additional names from Crane et al. 2004).
An updated phylogeny of Embryophytes based on the work by Novíkov & Barabaš-Krasni 2015 and Hao and Xue 2013 with plant taxon authors from Anderson, Anderson & Cleal 2007 and some additional clade names. Puttick et al./Nishiyama et al. are used for the basal clades.
Diversity
Non-vascular land plants
The non-vascular land plants, namely the mosses (Bryophyta), hornworts (Anthocerotophyta), and liverworts (Marchantiophyta), are relatively small plants, often confined to environments that are humid or at least seasonally moist. They are limited by their reliance on water needed to disperse their gametes; a few are truly aquatic. Most are tropical, but there are many arctic species. They may locally dominate the ground cover in tundra and Arctic–alpine habitats or the epiphyte flora in rain forest habitats.
They are usually studied together because of their many similarities. All three groups share a haploid-dominant (gametophyte) life cycle and unbranched sporophytes (the plant's diploid generation). These traits appear to be common to all early diverging lineages of non-vascular plants on the land. Their life-cycle is strongly dominated by the haploid gametophyte generation. The sporophyte remains small and dependent on the parent gametophyte for its entire brief life. All other living groups of land plants have a life cycle dominated by the diploid sporophyte generation. It is in the diploid sporophyte that vascular tissue develops. In some ways, the term "non-vascular" is a misnomer. Some mosses and liverworts do produce a special type of vascular tissue composed of complex water-conducting cells. However, this tissue differs from that of "vascular" plants in that these water-conducting cells are not lignified. It is unlikely that water-conducting cells in the mosses is homologous with the vascular tissue in "vascular" plants.
Like the vascular plants, they have differentiated stems, and although these are most often no more than a few centimeters tall, they provide mechanical support. Most have leaves, although these typically are one cell thick and lack veins. They lack true roots or any deep anchoring structures. Some species grow a filamentous network of horizontal stems, but these have a primary function of mechanical attachment rather than extraction of soil nutrients (Palaeos 2008).
Rise of vascular plants
During the Silurian and Devonian periods (around ), plants evolved which possessed true vascular tissue, including cells with walls strengthened by lignin (tracheids). Some extinct early plants appear to be between the grade of organization of bryophytes and that of true vascular plants (eutracheophytes). Genera such as Horneophyton have water-conducting tissue more like that of mosses, but a different life-cycle in which the sporophyte is branched and more developed than the gametophyte. Genera such as Rhynia have a similar life-cycle but have simple tracheids and so are a kind of vascular plant. It was assumed that the gametophyte dominant phase seen in bryophytes used to be the ancestral condition in terrestrial plants, and that the sporophyte dominant stage in vascular plants was a derived trait. However, the gametophyte and sporophyte stages were probably equally independent from each other, and that the mosses and vascular plants in that case are both derived, and have evolved in opposite directions.
During the Devonian period, vascular plants diversified and spread to many different land environments. In addition to vascular tissues which transport water throughout the body, tracheophytes have an outer layer or cuticle that resists drying out. The sporophyte is the dominant generation, and in modern species develops leaves, stems and roots, while the gametophyte remains very small.
Lycophytes and euphyllophytes
All the vascular plants which disperse through spores were once thought to be related (and were often grouped as 'ferns and allies'). However, recent research suggests that leaves evolved quite separately in two different lineages. The lycophytes or lycopodiophytes – modern clubmosses, spikemosses and quillworts – make up less than 1% of living vascular plants. They have small leaves, often called 'microphylls' or 'lycophylls', which are borne all along the stems in the clubmosses and spikemosses, and which effectively grow from the base, via an intercalary meristem. It is believed that microphylls evolved from outgrowths on stems, such as spines, which later acquired veins (vascular traces).
Although the living lycophytes are all relatively small and inconspicuous plants, more common in the moist tropics than in temperate regions, during the Carboniferous period tree-like lycophytes (such as Lepidodendron) formed huge forests that dominated the landscape.
The euphyllophytes, making up more than 99% of living vascular plant species, have large 'true' leaves (megaphylls), which effectively grow from the sides or the apex, via marginal or apical meristems. One theory is that megaphylls evolved from three-dimensional branching systems by first '' – flattening to produce a two dimensional branched structure – and then 'webbing' – tissue growing out between the flattened branches. Others have questioned whether megaphylls evolved in the same way in different groups.
Ferns and horsetails
The ferns and horsetails (the Polypodiophyta) form a clade; they use spores as their main method of dispersal. Traditionally, whisk ferns and horsetails were historically treated as distinct from 'true' ferns. Living whisk ferns and horsetails do not have the large leaves (megaphylls) which would be expected of euphyllophytes. This has probably resulted from reduction, as evidenced by early fossil horsetails, in which the leaves are broad with branching veins.
Ferns are a large and diverse group, with some 12,000 species. A stereotypical fern has broad, much divided leaves, which grow by unrolling.
Seed plants
Seed plants, which first appeared in the fossil record towards the end of the Paleozoic era, reproduce using desiccation-resistant capsules called seeds. Starting from a plant which disperses by spores, highly complex changes are needed to produce seeds. The sporophyte has two kinds of spore-forming organs or sporangia. One kind, the megasporangium, produces only a single large spore, a megaspore. This sporangium is surrounded by sheathing layers or integuments which form the seed coat. Within the seed coat, the megaspore develops into a tiny gametophyte, which in turn produces one or more egg cells. Before fertilization, the sporangium and its contents plus its coat is called an ovule; after fertilization a seed. In parallel to these developments, the other kind of sporangium, the microsporangium, produces microspores. A tiny gametophyte develops inside the wall of a microspore, producing a pollen grain. Pollen grains can be physically transferred between plants by the wind or animals, most commonly insects. Pollen grains can also transfer to an ovule of the same plant, either with the same flower or between two flowers of the same plant (self-fertilization). When a pollen grain reaches an ovule, it enters via a microscopic gap in the coat, the micropyle. The tiny gametophyte inside the pollen grain then produces sperm cells which move to the egg cell and fertilize it. Seed plants include two clades with living members, the gymnosperms and the angiosperms or flowering plants. In gymnosperms, the ovules or seeds are not further enclosed. In angiosperms, they are enclosed within the carpel. Angiosperms typically also have other, secondary structures, such as petals, which together form a flower.
Meiosis in sexual land plants provides a direct mechanism for repairing DNA in reproductive tissues. Sexual reproduction appears to be needed for maintaining long-term genomic integrity and only infrequent combinations of extrinsic and intrinsic factors allow for shifts to asexuality.
| Biology and health sciences | Embryophytes (Land plants) except vascular plants | Plants |
318869 | https://en.wikipedia.org/wiki/Iceland%20spar | Iceland spar | Iceland spar, formerly called Iceland crystal ( , ) and also called optical calcite, is a transparent variety of calcite, or crystallized calcium carbonate, originally brought from Iceland, and used in demonstrating the polarization of light.
Formation and composition
Iceland spar is a colourless, transparent variety of calcium carbonate (CaCO3). It crystallizes in the trigonal system, typically forming rhombohedral crystals. It has a Mohs hardness of 3 and exhibits double refraction, splitting a ray of light into two rays that travel at different speeds and directions.
Iceland spar forms in sedimentary environments, mainly limestone and dolomite rocks, but it also forms in hydrothermal veins and evaporite deposits. It precipitates from solutions rich in calcium and carbonate ions, influenced by temperature, pressure, and impurities.
The most common crystal structure of Iceland spar is rhombohedral, but other structures, such as scalenohedral or prismatic, can form depending on formation conditions. Iceland spar is primarily found in Iceland but can occur in different parts of the world with suitable geological conditions.
Characteristics and optical properties
Iceland spar is characterized by its large, readily cleavable crystals, easily divided into parallelepipeds. This feature makes it easily identifiable and workable. One of the most remarkable properties of Iceland spar is its birefringence, where the crystal's refractive index differs for light of different polarizations. When a ray of unpolarized light passes through the crystal, it is divided into two rays of mutually perpendicular polarization directed at various angles. This double refraction causes objects seen through the crystal to appear doubled.
Iceland spar possesses several optical properties other than double refraction and birefringence. It is highly transparent to visible light, allowing light to pass through with minimal absorption or scattering, which is ideal for optical applications requiring clarity. Iceland spar can produce vivid colours when viewed under polarized light due to its birefringent nature. This effect is known as the "Becke line" and can be used to determine a mineral's refractive index. Additionally, Iceland spar is optically active, meaning it can rotate the plane of polarization of light passing through it, a property resulting from its asymmetrical atomic arrangement. These optical properties contribute to the mineral's scientific use and aesthetic appeal.
Historical significance
Iceland spar holds historical importance in optics and the study of light. One of its most notable properties is its ability to exhibit double refraction. This phenomenon was first described by the Danish scientist Erasmus Bartholin in 1669, who observed it in a specimen of Iceland spar.
The study of double refraction in Iceland spar played a role in developing the wave theory of light. Scientists such as Christiaan Huygens, Isaac Newton, and Sir George Stokes studied this phenomenon and contributed to the understanding of light as a wave. Huygens, in particular, used double refraction to support his wave theory of light, in contrast to Newton's corpuscular theory. Augustin-Jean Fresnel published a complete explanation of double refraction in light polarization in the 1820s.
The understanding of double refraction in Iceland spar also led to the development of polarized light microscopy, which is used in various scientific fields to study the properties of materials. Iceland spar has been used historically in optical instruments like polarizing microscopes and navigation equipment.
Mining
Mines producing Iceland spar include many mines producing related calcite and aragonite. Iceland spar occurs in various locations worldwide, historically named after Iceland due to its abundance on the island. Other productive sources include China and the greater Sonoran Desert region, in Santa Eulalia, Chihuahua, Mexico, and New Mexico, United States. The clearest specimens, as well as the largest, have been from the Helgustaðir mine in Iceland.
Surveying tools and techniques are combined to reduce the risk and cost of exploration to identify deposits. Geological maps and remote sensing techniques, such as satellite imagery and aerial photography, are used for initial exploration and regional assessment to identify potential areas for further exploration. Geophysical surveys, including magnetometry, gravity surveys, and electromagnetic surveys, are then employed to detect anomalies indicating mineralization. Field mapping of surface geology and mineralogy also plays a role in identifying potential mineralization zones.
The mining process for Iceland spar varies based on the specific geological conditions of the deposit. Open-pit mining or quarrying is common for surface deposits. Once extracted, the calcite is processed to remove impurities, prepared for various applications, including optical instruments and jewelry, and used as a source of calcium carbonate in industries like construction and agriculture.
Environmental issues
Some potential environmental issues associated with Iceland spar mining include habitat destruction, water pollution, air pollution, soil degradation, and visual impact. Mining activities can destroy natural habitats, mainly if the mining site is located in ecologically sensitive areas, leading to the loss of biodiversity and disrupting local ecosystems. Water sources can be contaminated through the discharge of chemicals used in the extraction and processing of minerals, impacting aquatic life and water quality. Mining activities can also lead to soil erosion and degradation, mainly if proper land reclamation measures are not implemented after mining ceases. Open-pit mining operations can have a significant visual impact on the landscape, altering the natural scenery of an area. These measures may include erosion control, environmentally friendly mining techniques, and the reclamation of mined areas to restore them to a natural state.
Health concerns
Mining, including Iceland spar mining, poses various health risks to workers and nearby communities. Some key health concerns associated with mining activities include respiratory issues, noise-induced hearing loss, chemical exposure, musculoskeletal disorders, injuries and accidents, and mental health issues. Dust generated during mining operations can contain harmful particles, leading to respiratory problems. The high noise levels generated by mining activities can cause hearing loss over time if proper protective measures are not in place. Miners may also be exposed to harmful chemicals used in the extraction and processing of minerals, which can cause various health issues. The physical demands of mining work, such as heavy lifting and repetitive motions, can result in musculoskeletal disorders. Injuries and accidents are also common risks in mining, including falls, equipment-related incidents, and mine collapses. The demanding nature of mining work, along with long hours and isolation, can contribute to mental health issues such as stress, anxiety, and depression. Mining companies must implement health and safety measures to mitigate these risks to protect workers and nearby communities, including personal protective equipment, dust control measures, and health and safety training. Regularly monitoring air quality, noise levels, and other potential hazards is essential to ensure a safe working environment.
Uses
Iceland spar has been historically used in telecommunications due to its unique optical properties. One of its key features, birefringence, made it worthwhile in early optical technologies, such as developing optical instruments like polarizing microscopes and constructing optical rangefinders and gunsights.
While uncommon, Iceland spar has historically been used in navigation as a polarizing filter to determine the sun's direction on overcast days. It has been speculated that the sunstone (, a different mineral from the gem-quality sunstone) mentioned in medieval Icelandic texts, such as Rauðúlfs þáttr, was Iceland spar, and that Vikings used its light-polarizing property to tell the direction of the sun on cloudy days for navigational purposes. The polarization of sunlight in the Arctic can be detected, and the direction of the sun identified to within a few degrees in both cloudy and twilight conditions using the sunstone and the naked eye. The process involves moving the stone across the visual field to reveal a yellow entoptic pattern on the fovea of the eye, probably Haidinger's brush. The recovery of an Iceland spar sunstone from a ship of the Elizabethan era that sank in 1592 off Alderney suggests that this navigational technology may have persisted after the invention of the magnetic compass.
William Nicol (1770–1851) invented the first polarizing prism, using Iceland spar to create his Nicol prism.
Modern applications
Despite being historically significant, Iceland spar still holds an essential place in modern applications. Due to its optical properties, Iceland spar is still used in instruments like polarizing microscopes, lenses, and filters. Iceland spar is also used in optical instruments for geological and biological microscopy as its birefringence helps to reveal material structure. It is also a practical tool used in education and research to demonstrate optical principles. Though its applications are less widespread than in the past, Iceland spar continues to contribute to various scientific and technological endeavours.
As a type of calcite, Iceland spar can be used in construction as a building material in cement and concrete. Its high purity and brightness make it an ideal filler in paints and coatings. In metallurgy, calcite acts as a flux to lower the melting point of metals during smelting and refining. Additionally, it is used in agriculture as a soil conditioner and neutralizer to adjust soil pH levels and improve crop yields. Calcite also contributes to environmental remediation efforts, treating acidic water and soil by neutralizing acidity and removing heavy metals.
Geological significance
Due to Iceland spar typically forming in sedimentary environments, particularly limestone and dolomite rocks, its formation is closely tied to these carbonate rocks' deposition and diagenesis (compaction and cementation). Studying Iceland spar can provide valuable information about past environmental conditions, such as the presence of ancient seas and marine life, as carbonate rocks like limestone often form in marine environments. The presence of Iceland spar can also indicate the presence of hydrothermal activity, as calcite can form in hydrothermal veins.
Conservation and protection
Due to their scientific and historical significance, conservation efforts related to Iceland spar primarily focus on preserving specimens and mining sites. One of the challenges in preserving Iceland spar specimens is the risk of damage during extraction, handling, and storage. Mining sites that yield high-quality Iceland spar specimens are also of interest for conservation. These sites may be designated protected areas to prevent overexploitation.
Cultural impact
The Thomas Pynchon novel Against the Day uses the doubling effect of Iceland spar as a theme.
| Physical sciences | Minerals | Earth science |
318895 | https://en.wikipedia.org/wiki/Deodorant | Deodorant | A deodorant is a substance applied to the body to prevent or mask body odor caused by bacterial breakdown of perspiration, for example in the armpits, groin, or feet. A subclass of deodorants, called antiperspirants, prevents sweating itself, typically by blocking sweat glands. Antiperspirants are used on a wider range of body parts, at any place where sweat would be inconvenient or unsafe, since unwanted sweating can interfere with comfort, vision, and grip (due to slipping). Other types of deodorant allow sweating but prevent bacterial action on sweat, since human sweat only has a noticeable smell when it is decomposed by bacteria.
In the United States, the Food and Drug Administration classifies and regulates most deodorants as cosmetics, but classifies antiperspirants as over-the-counter drugs.
The first commercial deodorant, Mum, was introduced and patented in the late nineteenth century by an inventor in Philadelphia, Pennsylvania, Edna Murphey. The product was briefly withdrawn from the market in the US. The modern formulation of the antiperspirant was patented by Jules Montenier on January 28, 1941. This formulation was first found in "Stopette" deodorant spray, which Time magazine called "the best-selling deodorant of the early 1950s".
Use of deodorant with aluminium compounds has been suspected of being linked to breast cancer, but research has not proven any such link.
Overview
The human body produces perspiration (sweat) via two types of sweat gland: eccrine sweat glands which cover much of the skin and produce watery odourless sweat, and apocrine sweat glands in the armpits and groin, which produce a more oily "heavy" sweat containing a proportion of waste proteins, fatty acids and carbohydrates, that can be metabolized by bacteria to produce compounds that cause body odor. In addition, the vagina produces secretions which are not a form of sweat but may be undesired and also masked with deodorants.
Human perspiration of all types is largely odorless until its organic components are fermented by bacteria that thrive in hot, humid environments. The human underarm is among the most consistently warm areas on the surface of the human body, and sweat glands readily provide moisture containing a fraction of organic matter, which when excreted, has a vital cooling effect. When adult armpits are washed with alkaline pH soap, the skin loses its protective acid mantle (pH 4.5–6), raising the skin pH and disrupting the skin barrier. Many bacteria are adapted to the slightly alkaline environment within the human body, so they can thrive within this elevated pH environment. This makes the skin more than usually susceptible to bacterial colonization. Bacteria on the skin feed on the waste proteins and fatty acids in the sweat from the apocrine glands and on dead skin and hair cells, releasing trans-3-methyl-2-hexenoic acid in their waste, which is the primary cause of body odor.
Underarm hair wicks the moisture away from the skin and aids in keeping the skin dry enough to prevent or diminish bacterial colonization. The hair is less susceptible to bacterial growth and therefore reduces bacterial odor. The apocrine sweat glands are inactive until puberty, which is why body odor often only becomes noticeable at that time.
Deodorant products work in one of two ways – by preventing sweat from occurring, or by allowing it to occur but preventing bacterial activity that decomposes sweat on the skin.
History
Modern deodorants
In 1888, the first modern commercial deodorant, Mum, was developed and patented by a U.S. inventor in Philadelphia, Pennsylvania, Edna Murphey; the small company was bought by Bristol-Myers in 1931. In the late 1940s, Helen Barnett Diserens developed an underarm applicator based on the newly invented ball-point pen. In 1952, the company began marketing the product under the name Ban Roll-On. The product was briefly withdrawn from the market in the U.S., but it is once again available at retailers in the U.S. under the brand Ban. In the UK it is sold under the names Mum Solid and Mum Pump Spray. Chattem acquired the Ban deodorant brand in 1998 and subsequently sold it to Kao Corporation in 2000.
In 1903, the first commercial antiperspirant was Everdry. The modern formulation of the antiperspirant was patented by Jules Montenier on January 28, 1941. This patent addressed the problem of the excessive acidity of aluminum chloride and its excessive irritation of the skin, by combining it with a soluble nitrile or a similar compound. This formulation was first found in "Stopette" deodorant spray, which Time magazine called "the best-selling deodorant of the early 1950s". "Stopette" gained its prominence as the first and long-time sponsor of the game show What's My Line?; it was later eclipsed by many other brands once the 1941 patent expired.
Between 1942 and 1957, the market for deodorants increased 600 times to become a $70 million market. Deodorants were originally marketed primarily to women, but by 1957 the market had expanded to male users, and estimates were that 50% of men were using deodorants by that date. The Ban Roll-On product led the market in sales.
In the early 1960s, the first aerosol antiperspirant in the marketplace was Gillette's Right Guard, whose brand was later sold to Henkel in 2006. Aerosols were popular because they let the user dispense a spray without coming in contact with the underarm area. By the late 1960s, half of all the antiperspirants sold in the U.S. were aerosols, and continued to grow in all sales to 82% by the early 1970s. However, the late 1970s saw two developments which greatly reduced the popularity of these products. First, in 1977 the U.S. Food and Drug Administration banned the active ingredient used in aerosols, aluminium zirconium chemicals, due to safety concerns over long term inhalation. Second, the U.S. Environmental Protection Agency limited the use of chlorofluorocarbon (CFC) propellants used in aerosols due to awareness that these gases can contribute to depleting the ozone layer. As the popularity of aerosols slowly decreased, stick antiperspirants became more popular.
Classification
Deodorant
In the United States, deodorants are classified and regulated as cosmetics by the U.S. Food and Drug Administration (FDA) and are designed to eliminate odor. Deodorants are often alcohol-based. Alcohol initially stimulates sweating but may also temporarily kill bacteria. Other active ingredients in deodorants include sodium stearate, sodium chloride, and stearyl alcohol. Deodorants can be formulated with other, more persistent antimicrobials such as triclosan that slow bacterial growth or with metal chelant compounds such as EDTA. Deodorants may contain perfume fragrances or natural essential oils intended to mask the odor of perspiration. Some of the first patented deodorants used zinc oxide, acids, ammonium chloride, sodium bicarbonate, and formaldehyde (which is now known as a carcinogen), and some of these ingredients were messy, irritating to the skin.
Over-the-counter products, often labeled as "natural deodorant crystal", contain the chemical rock crystals potassium alum or ammonium alum, which prevents bacterial action on sweat. These have gained popularity as an alternative health product, in spite of concerns about possible risks related to aluminum (see below – all alum salts contain aluminum in the form of aluminum sulphate salts) and contact dermatitis.
Vaginal deodorant, in the form of sprays, suppositories, and wipes, is often used by women to mask vaginal secretions. Vaginal deodorants can sometimes cause dermatitis.
Deodorant antiperspirant
In the United States, deodorants combined with antiperspirant agents are classified as drugs by the FDA. Antiperspirants attempt to stop or significantly reduce perspiration and thus reduce the moist climate in which bacteria thrive. Aluminium chloride, aluminium chlorohydrate, and aluminium-zirconium compounds, most notably aluminium zirconium tetrachlorohydrex gly are frequently used in antiperspirants. Aluminium chlorohydrate and aluminium-zirconium tetrachlorohydrate gly are the most frequent active ingredients in commercial antiperspirants. Aluminium-based complexes react with the electrolytes in the sweat to form a gel plug in the duct of the sweat gland. The plugs prevent the gland from excreting liquid and are removed over time by the natural sloughing of the skin. The metal salts work in another way to prevent sweat from reaching the surface of the skin: the aluminium salts interact with the keratin fibrils in the sweat ducts and form a physical plug that prevents sweat from reaching the skin's surface. Aluminium salts also have a slight astringent effect on the pores; causing them to contract, further preventing sweat from reaching the surface of the skin. The blockage of a large number of sweat glands reduces the amount of sweat produced in the underarms, though this may vary from person to person. Methenamine in the form of cream or spray is effective in the treatment of excessive sweating and attendant odor. Antiperspirants are usually best applied before bed.
Product formulations and formats
Formulations
Common and historical formulations for deodorants include the following active ingredients:
Aluminum salts (aluminum chlorohydrate, aluminum zirconium tetrachlorohydrex gly, and others) – used as the basis for almost all non-prescription (everyday) antiperspirants. The aluminum reacts within the sweat gland to form a colloid which physically prevents sweating.
Alum (typically potassium alum or ammonium alum, also described as "rock alum", or "rock crystal", or "natural deodorant"). Alum is a natural crystalline product widely used both historically and in modern times as a deodorant, because it inhibits bacterial action. The word 'alum' is a historical term for aluminum sulfate salts, therefore all alum products will contain aluminum, albeit in a different chemical form from antiperspirants.
Bactericidal products such as triclosan (TCS), octenidine dihydrochloride, and parabens kill bacteria on the skin.
Alcohols and related compounds such as propylene glycol – these products can have both drying and bactericidal effects.
Methenamine (hexamethylenetetramine, also known as hexamine or urotropin) is a powerful antiperspirant, often used for severe sweat-related issues, as well as prevention of sweating within the sockets of prosthetic devices used by amputees.
Acidifiers and pH neutral products – deodorants that prevent bacterial action by enhancing (or at least, not depleting) the skin's natural slight acidity, known as the acid mantle, which naturally reduces bacterial action but can be compromised by typically alkaline soaps and skin products.
Masking scents – other strong or overriding scents of a pleasing type may be used, used to mask bodily odors. Typically these are strongly smelling plant extracts or synthetic aromas.
Activated charcoal and other products capable of absorbing sweat and/or smell. Although charcoal most often has a black color, the activated charcoal used in deodorants may be a very light color for aesthetic reasons.
Less commonly used, products such as milk of magnesia (a thick liquid suspension of magnesium hydroxide) are sometimes used as deodorants. Many milk of magnesia products contain small amounts of sodium hypochlorite (bleach) at very low levels that are safe for ingestion and skin application. Sodium hypochlorite is a powerful bactericide, and it is possible that its presence in a product that can dry onto the skin, may explain this use as a deodorant. (Safety info: bleach is caustic and extremely poisonous, and can be lethal, in higher concentrations.)
Formats
Deodorants and antiperspirants come in many forms. What is commonly used varies in different countries. In Europe, aerosol sprays are popular, as are cream and roll-on forms. In North America, solid or gel forms are dominant.
Health effects
After using a deodorant containing zirconium, the skin may develop an allergic, axillary granuloma response. Antiperspirants with propylene glycol, when applied to the axillae, can cause irritation and may promote sensitization to other ingredients in the antiperspirant. Deodorant crystals containing synthetically made potassium alum were found to be a weak irritant to the skin. Unscented deodorant is available for those with sensitive skin. Frequent use of deodorants was associated with blood concentrations of the synthetic musk galaxolide.
Aluminum
Many deodorants and antiperspirants contain aluminium in the form of aluminium salts such as aluminium chlorohydrate.
The US Food and Drug Administration, in a 2003 paper discussing deodorant safety, concluded that "despite many investigators looking at this issue, the agency does not find data from topical and inhalation chronic exposure animal and human studies submitted to date sufficient to change the monograph status of aluminum containing antiperspirants", therefore allowing their use and stating they will keep monitoring the scientific literature. Members of the Scientific Committee on Consumer Safety (SCCS) of the European Commission concluded similarly in 2015, that "due to the lack of adequate data on dermal penetration to estimate the internal dose of aluminium following cosmetic uses, risk assessment cannot be performed." In the light of new data in 2020 the SCCS considered aluminium compounds safe up to 6.25% in non-spray deodorants or non-spray antiperspirants and 10.60% in spray deodorants or spray antiperspirants.
Myths and claims related to aluminium compounds in deodorants
Common myths and marketing claims for aluminium in deodorants (including aluminum in alum products) include claims:
That aluminium in deodorants applied to the skin is a risk factor for some cancers (notably breast cancer) and some forms of dementia
That aluminium in antiperspirants can enter the body (possibly through shaving cuts)
That aluminium in alum "natural deodorant" products is "safer" because it is "too large" to enter the body
Of note, the parts of the body which are commonly shaved and also commonly treated with deodorants, such as the armpits, contain substantial deposits of subcutaneous fat. Shaving cuts would be extremely unlikely to penetrate sufficiently beyond the very outer layers of the skin, for much if any product to enter the bloodstream.
Alzheimer's disease
A 2014 review of 469 peer-reviewed studies examining the effect of exposure to aluminum products concluded "that health risks posed by exposure to inorganic depend on its physical and chemical forms and that the response varies with route of administration, magnitude, duration and frequency of exposure. These results support previous conclusions that there is little evidence that exposure to metallic Al, the Al oxides or its salts increases risk for Alzheimer's disease, genetic damage or cancer".
Breast cancer
The claim that breast cancer is believed to be linked with deodorant use has been widely circulated and appears to originate from a spam email sent in 1999; however, there is no evidence to support the existence of such a link. The myth circulates in two forms:
Antiperspirants block the "purging" of toxins which build up in the body and cause breast cancer: As sweat glands simply do not have this function, the claim is scientifically implausible. Perspiration from the eccrine sweat glands is 99% water, with some salt (sodium chloride) and only trace amounts of lactic acid (almost entirely processed in the liver), urea (almost entirely excreted by the kidneys), and only very small amounts of all other components. Perspiration from the apocrine sweat glands (those in the armpits and groin, which are more responsible for body odor) also include waste proteins, carbohydrates, and fatty acids which would otherwise be processed by other organs such as the liver.
It is possible that there has been confusion between sweat glands, and the lymph nodes deep within the armpits which form part of the immune system and help filter toxins, but if so, there is no evidence at all of such "blocking" of lymph nodes, nor any scientifically plausible route by which this could result from deodorant use.
Aluminum in antiperspirants can enter the body (possibly through cuts) and cause breast cancer: There is no current evidence to support this claim, nor any convincing evidence that it is true. A fact often cited to back up this claim is that more breast cancers occur in the part of the breast near the armpits. However, breast tissue is not evenly spread out, and the part of the breast near the armpit (the Tail of Spence) simply contains much more breast tissue than the other quadrants, making it much more likely that any cancer would occur in that location. See above for current scientific knowledge regarding aluminum in deodorants.
The National Cancer Institute states that "no scientific evidence links the use of these products to the development of breast cancer" and that "no clear evidence that the use of aluminum-containing underarm antiperspirants or cosmetics increases the risk of breast cancer", but also concludes that studies of antiperspirants and deodorants and breast cancer have provided conflicting results, additional research would be needed to determine whether a relationship exists".
Another constituent of deodorant products that has given cause for concern are parabens, a chemical additive. However parabens do not cause cancer.
Kidney dysfunction
The FDA has "acknowledge[d] that small amounts of aluminum can be absorbed from the gastrointestinal tract and through the skin", leading to a warning "that people with kidney disease may not be aware that the daily use of antiperspirant drug products containing aluminum may put them at a higher risk because of exposure to aluminum in the product." The agency warns people with kidney dysfunction to consult a doctor before using antiperspirants containing aluminum.
Aerosol burns and frostbite
If aerosol deodorant is held close to the skin for long enough, it can cause an aerosol burn—a form of frostbite. In controlled tests, spray deodorants have been shown to cause temperature drops of over 60 °C in a short period of time.
Clothing
Aluminium zirconium tetrachlorohydrex gly, a common antiperspirant, can react with sweat to create yellow stains on clothing. Underarm liners are an antiperspirant alternative that does not leave stains.
| Biology and health sciences | Hygiene products | Health |
318980 | https://en.wikipedia.org/wiki/Endorheic%20basin | Endorheic basin | An endorheic basin ( ; also endoreic basin and endorreic basin) is a drainage basin that normally retains water and allows no outflow to other external bodies of water (e.g. rivers and oceans); instead, the water drainage flows into permanent and seasonal lakes and swamps that equilibrate through evaporation. Endorheic basins are also called closed basins, terminal basins, and internal drainage systems.
Endorheic regions contrast with open lakes (exorheic regions), where surface waters eventually drain into the ocean. In general, water basins with subsurface outflows that lead to the ocean are not considered endorheic; but cryptorheic. Endorheic basins constitute local base levels, defining a limit of the erosion and deposition processes of nearby areas. Endorheic water bodies include the Caspian Sea, which is the world's largest inland body of water.
Etymology
The term endorheic derives from the French word , which combines ( 'within') and 'flow'.
Endorheic lakes
Endorheic lakes (terminal lakes) are bodies of water that do not flow into an ocean or a sea. Most of the water that falls to Earth percolates into the oceans and the seas by way of a network of rivers, lakes, and wetlands. Analogous to endorheic lakes is the class of bodies of water located in closed watersheds (endorheic watersheds) where the local topography prevents the drainage of water into the oceans and the seas. These endorheic watersheds (containing water in rivers or lakes that form a balance of surface inflows, evaporation and seepage) are often called sinks.
Endorheic lakes are typically located in the interior of a landmass, far from an ocean, and in areas of relatively low rainfall. Their watersheds are often confined by natural geologic land formations such as a mountain range, cutting off water egress to the ocean. The inland water flows into dry watersheds where the water evaporates, leaving a high concentration of minerals and other inflow erosion products. Over time this input of erosion products can cause the endorheic lake to become relatively saline (a "salt lake"). Since the main outflow pathways of these lakes are chiefly through evaporation and seepage, endorheic lakes are usually more sensitive to environmental pollutant inputs than water bodies that have access to oceans, as pollution can be trapped in them and accumulate over time.
Occurrence
Endorheic regions can occur in any climate but are most commonly found in desert locations. This reflects the balance between tectonic subsidence and rates of evaporation and sedimentation. Where the basin floor is dropping more rapidly than water and sediments can accumulate, any lake in the basin will remain below the sill level (the level at which water can find a path out of the basin). Low rainfall or rapid evaporation in the watershed favor this case. In areas where rainfall is higher, riparian erosion will generally carve drainage channels (particularly in times of flood), or cause the water level in the terminal lake to rise until it finds an outlet, breaking the enclosed endorheic hydrological system's geographical barrier and opening it to the surrounding terrain. The Black Sea was likely such a lake, having once been an independent hydrological system before the Mediterranean Sea broke through the terrain separating the two. Lake Bonneville was another such lake, overflowing its basin in the Bonneville flood. The Malheur/Harney lake system in Oregon is normally cut off from drainage to the ocean, but has an outflow channel to the Malheur River. This is presently dry, but may have flowed as recently as 1,000 years ago.
Examples of relatively humid regions in endorheic basins often exist at high elevation. These regions tend to be marshy and are subject to substantial flooding in wet years. The area containing Mexico City is one such case, with annual precipitation of and characterized by waterlogged soils that require draining.
Endorheic regions tend to be far inland with their boundaries defined by mountains or other geological features that block their access to oceans. Since the inflowing water can evacuate only through seepage or evaporation, dried minerals or other products collect in the basin, eventually making the water saline and also making the basin vulnerable to pollution. Continents vary in their concentration of endorheic regions due to conditions of geography and climate. Australia has the highest percentage of endorheic regions at 21 per cent while North America has the least at five per cent. Approximately 18 per cent of the Earth's land drains to endorheic lakes or seas, the largest of these land areas being the interior of Asia.
In deserts, water inflow is low and loss to solar evaporation high, drastically reducing the formation of complete drainage systems. In the extreme case, where there is no discernible drainage system, the basin is described as arheic. Closed water flow areas often lead to the concentration of salts and other minerals in the basin. Minerals leached from the surrounding rocks are deposited in the basin, and left behind when the water evaporates. Thus endorheic basins often contain extensive salt pans (also called salt flats, salt lakes, alkali flats, dry lake beds, or playas). These areas tend to be large, flat hardened surfaces and are sometimes used for aviation runways, or land speed record attempts, because of their extensive areas of perfectly level terrain.
Both permanent and seasonal endorheic lakes can form in endorheic basins. Some endorheic basins are essentially stable because climate change has reduced precipitation to the degree that a lake no longer forms. Even most permanent endorheic lakes change size and shape dramatically over time, often becoming much smaller or breaking into several smaller parts during the dry season. As humans have expanded into previously uninhabitable desert areas, the river systems that feed many endorheic lakes have been altered by the construction of dams and aqueducts. As a result, many endorheic lakes in developed or developing countries have contracted dramatically, resulting in increased salinity, higher concentrations of pollutants, and the disruption of ecosystems.
Even within exorheic basins, there can be "non-contributing", low-lying areas that trap runoff and prevent it from contributing to flows downstream during years of average or below-average runoff. In flat river basins, non-contributing areas can be a large fraction of the river basin, e.g. Lake Winnipeg's basin. A lake may be endorheic during dry years and can overflow its basin during wet years, e.g., the former Tulare Lake.
Because the Earth's climate has recently been through a warming and drying phase with the end of the Ice Ages, many endorheic areas such as Death Valley that are now dry deserts were large lakes relatively recently. During the last ice age, the Sahara may have contained lakes larger than any now existing.
Climate change coupled with the mismanagement of water in these endorheic regions has led to devastating losses in ecosystem services and toxic surges of pollutants. The desiccation of saline lakes produces fine dust particles that impair agriculture productivity and harm human health. Anthropogenic activity has also caused a redistribution of water from these hydrologically landlocked basins such that endorheic water loss has contributed to sea level rise, and it is estimated that most of the terrestrial water lost ends up in the ocean. In regions such as Central Asia, where people depend on endorheic basins and other surface water sources to satisfy their water needs, human activity greatly impacts the availability of that water.
Notable endorheic basins and lakes
Africa
Large endorheic regions in Africa are located in the Sahara Desert, the Sahel, the Kalahari Desert, and the East African Rift:
Chad Basin, in the northern centre of Africa. It covers an area of approximately 2.434 million km2.
Qattara Depression, in Egypt.
Chott Melrhir, in Algeria.
Chott el Djerid, in Tunisia.
The Okavango River, in the Kalahari Desert, is part of an endorheic basin region, the Okavango Basin, that also includes the Okavango Delta, Lake Ngami, the Nata River, and a number of salt pans such as Makgadikgadi Pan.
Etosha Pan in Namibia's Etosha National Park.
Turkana Basin, in Kenya, whose basin includes the Omo River of Ethiopia.
Lake Chilwa, in Malawi.
Afar Depression, in Eritrea, Ethiopia, and Djibouti, which contains the Awash River
Some Rift Valley lakes, such as Lake Abijatta, Lake Chew Bahir, Lake Shala, Lake Chamo, and Lake Awasa.
Lake Mweru Wantipa, in Zambia.
Lake Magadi, in Kenya.
Lake Rukwa, in Tanzania.
Antarctica
Endorheic lakes exist in Antarctica's McMurdo Dry Valleys, Victoria Land, the largest ice-free area.
Don Juan Pond in Wright Valley is fed by groundwater from a rock glacier and remains unfrozen throughout the year.
Lake Vanda in Wright Valley has a perennial ice cover, the edges of which melt in the summer, allowing flow from the longest river in Antarctica, the Onyx River. The lake is over 70 m deep and is hypersaline.
Lake Bonney is in Taylor Valley and has a perennial ice cover and two lobes separated by the Bonney Riegel. Glacial melt and discharge from Blood Falls feed the lake. Its unique glacial history has resulted in hypersaline brine in the bottom waters and fresh water at the surface.
Lake Hoare, in Taylor Valley, is the freshest of the Dry Valley lakes, receiving its melt almost exclusively from the Canada Glacier. The lake has an ice cover and forms a moat during the Austral summer.
Lake Fryxell is adjacent to the Ross Sea in Taylor Valley. The lake has an ice cover and receives its water from numerous glacial meltwater streams for approximately six weeks out of the year. Its salinity increases with depth.
Asia
Much of Western and Central Asia is a giant endorheic region made up of a number of contiguous closed basins. The region contains several basins and terminal lakes, including:
The Caspian Sea, the largest lake on Earth. A large part of western Russia, drained by the Volga River, is part of the Caspian basin.
Lake Urmia in Western Azerbaijan Province of Iran.
The Aral Sea, whose tributary rivers have been diverted, leading to a dramatic shrinkage of the lake. The resulting ecological disaster has brought the plight of internal drainage basins to public attention.
Lake Balkhash, in Kazakhstan.
Issyk-Kul Lake and Chatyr-Kul Lake in Kyrgyzstan.
Lop Lake, in the Tarim Basin of China's Xinjiang Uygur Autonomous Region.
The Dzungarian Basin in Xinjiang, separated from the Tarim Basin by the Tian Shan. The most notable terminal lake in the basin is the Manas Lake.
The Central Asian Internal Drainage Basin, in southern and western Mongolia, contains a series of closed drainage basins, such as the Khyargas Nuur basin, the Uvs Nuur basin, which includes Üüreg Lake, and the Pu-Lun-To River Basin.
Qaidam Basin, in Qinghai Province, China, as well as nearby Qinghai Lake.
Sistan Basin covering areas of Iran and Afghanistan
Pangong Tso and Aksai Chin Lake on the China-India border
Many small lakes and rivers of the Iranian Plateau, including Gavkhouni marshes and Namak Lake.
Other endorheic lakes and basins in Asia include:
The Dead Sea, the lowest surface point on Earth and one of its saltiest bodies of water lies between Israel and Jordan.
Sambhar Lake, in Rajasthan, north-western India
Lake Van in eastern Turkey
Sabkhat al-Jabbul, extensive salt flats and a lake in Syria.
Solar Lake, Sinai, near the Israeli-Egypt border.
Lake Tuz, in Turkey, in south part of Central Anatolia Region.
Sawa lake in Iraq, in Muthanna Governorate.
Australia
Australia, being very dry and having exceedingly low runoff ratios due to its ancient soils, has many endorheic drainages. The most important are:
Lake Eyre basin, which drains into the highly variable Lake Eyre and includes Lake Frome.
Lake Torrens, usually an endorheic lake to the west of the Flinders Ranges in South Australia, that flows to the sea after extreme rainfall events.
Lake Corangamite, a highly saline crater lake in western Victoria.
Lake George, formerly connected to the Murray-Darling Basin
Europe
Though a large portion of Europe drains to the endorheic Caspian Sea, Europe's wet climate means it contains relatively few terminal lakes itself: any such basin is likely to continue to fill until it reaches an overflow level connecting it with an outlet or erodes the barrier blocking its exit.
There are some seemingly endorheic lakes, but they are cryptorheic, being drained either through manmade canals, via karstic phenomena, or other subsurface seepage.
Lake Neusiedl, in Austria and Hungary.
Lake Trasimeno, in Italy.
Fucine Lake, in Italy. Now drained.
Lake Velence, in Hungary.
Lake Prespa, between Albania, Greece and North Macedonia.
Rahasane Turlough, the largest turlough in Ireland.
Laacher See, in Germany.
The Lasithi Plateau in Crete, Greece, is a high endorheic plateau.
A few minor true endorheic lakes exist in Spain (e.g. Laguna de Gallocanta, Estany de Banyoles), Italy, Cyprus (Larnaca and Akrotiri salt lakes) and Greece.
North and Central America
The Great Basin is North America's largest and the world's ninth largest endorheic basin, covering nearly all of Nevada, much of Oregon and Utah, and portions of California, Idaho, and Wyoming. Notable enclosed basins include Death Valley, the hottest location on Earth; the Black Rock Desert and Bonneville Salt Flats, location of many of the new vehicle land speed records set since the 1930s; the Great Salt Lake, remnant of Lake Bonneville; and the Salton Sea.
The Valley of Mexico. In Pre-Columbian times, the Valley was substantially covered with five lakes, including Lake Texcoco, Lake Xochimilco, and Lake Chalco.
Guzmán Basin, in northern Mexico and the southwestern United States. The Mimbres River of New Mexico drains into this basin.
Lago Atitlán, a volcanic caldera lake in the highlands of Guatemala. It is cryptorheic.
Lago Coatepeque, El Salvador.
Bolsón de Mapimí, in northern Mexico.
Willcox Playa of southern Arizona.
Tulare Lake in the San Joaquin Valley in Central California, fed by the Kaweah and Tule Rivers plus southern distributaries of the Kings. Historically, it would drain into the San Joaquin River in very wet years. Agricultural development and irrigation diversions have left the lake dry.
Buena Vista Lake at the southmost end of the San Joaquin Valley in Southern California, fed by the Kern River. Historically, it would drain into Tulare Lake and the San Joaquin River in exceptionally wet years. Agricultural development and irrigation diversions have left the lake dry.
Crater Lake, in Oregon, a cryptorheic lake with subsurface drainage to the Wood River. It is filled directly by rain and snow and has very little mineral or salt buildup.
The Great Divide Basin in Wyoming, a small endorheic basin that straddles the Continental Divide of the Americas.
Devils Lake, in North Dakota.
Devil's Lake, in Wisconsin, cryptorheic.
Tule Lake and the Lost River basin in California and Oregon.
Little Manitou Lake in Saskatchewan.
Old Wives Lake, on the Laurentian Divide in Saskatchewan.
Quill Lakes, in Saskatchewan.
Pakowki Lake, on the Laurentian Divide in Alberta.
Paynes Prairie, in Florida. Since 1927, it has been drained by canal to the Atlantic Ocean via the River Styx.
Spotted Lake, Osoyoos, British Columbia, Canada.
Several lakes on the western Chilcotin Plateau sit on the divide between the Fraser River drainage to the east and the Homathko drainage to the west. Such examples include Choelquoit Lake, Eagle Lake, and Martin lake.
Frame Lake in Yellowknife, capital of Canada's Northwest Territories.
New Mexico has several desert endorheic basins, including:
The Tularosa Basin, a rift valley.
Zuñi Salt Lake, a maar.
The Mimbres River Basin, in Grant County.
The San Agustin Basin, in Catron and Socorro Counties.
Lago Enriquillo on the island of Hispaniola in the Caribbean Sea.
Many small lakes and ponds in North Dakota and the Northern Great Plains are endorheic, and some have salt encrustations along their shores.
South America
Laguna del Carbón, in Gran Bajo de San Julián, Argentina – the lowest point in the Western and Southern hemispheres
Lake Mar Chiquita in Argentina.
The Altiplano includes a number of closed basins such as the Salar de Coipasa, and Titicaca–Poopó system.
Lake Valencia, in Venezuela.
Salar de Atacama, in the Atacama Desert, Chile.
Ancient
Some of Earth's ancient endorheic systems and lakes include:
The Black Sea, until its merger with the Mediterranean.
The Mediterranean Sea itself and all its tributary basins, during its Messinian desiccation (approximately five million years ago) as it became disconnected from the Atlantic Ocean.
The Orcadian Basin in Scotland during the Devonian period. Now identifiable as lacustrine sediments buried around and off the coast.
Lake Tanganyika in Africa. Currently high enough to connect to rivers entering the sea.
Lake Lahontan in North America.
Lake Bonneville in North America. The basin was not always endorheic; at times, it overflowed through Red Rock Pass to the Snake River and the sea.
Lake Chewaucan in North America.
Tularosa Basin and Lake Cabeza de Vaca in North America. The basin was formerly much larger than it is today, including the ancestral Rio Grande north of Texas, which fed a large lake area.
Ebro and Duero basins, draining most of northern Spain during the Neogene and perhaps Pliocene. Climate change and erosion of the Catalan coastal mountains, as well as the deposition of alluvium in the terminal lake, allowed the Ebro basin to overflow into the sea during the middle-to-late Miocene.
| Physical sciences | Hydrology | Earth science |
319122 | https://en.wikipedia.org/wiki/Pinwheel%20Galaxy | Pinwheel Galaxy | The Pinwheel Galaxy (also known as Messier 101, M101 or NGC 5457) is a face-on, unbarred, and counterclockwise spiral galaxy located from Earth in the constellation Ursa Major. It was discovered by Pierre Méchain in 1781 and was communicated that year to Charles Messier, who verified its position for inclusion in the Messier Catalogue as one of its final entries.
On February 28, 2006, NASA and the European Space Agency released a very detailed image of the Pinwheel Galaxy, which was the largest and most detailed image of a galaxy by Hubble Space Telescope at the time. The image was composed of 51 individual exposures, plus some extra ground-based photos.
Discovery
Pierre Méchain, the discoverer of the galaxy, described it as a "nebula without star, very obscure and pretty large, 6' to 7' in diameter, between the left hand of Bootes and the tail of the great Bear. It is difficult to distinguish when one lits the [grating] wires."
William Herschel wrote in 1784 that the galaxy was one of several which "...in my 7-, 10-, and 20-feet [focal length] reflectors shewed a mottled kind of nebulosity, which I shall call resolvable; so that I expect my present telescope will, perhaps, render the stars visible of which I suppose them to be composed."
Lord Rosse observed the galaxy in his 72-inch-diameter Newtonian reflector during the second half of the 19th century. He was the first to make extensive note of the spiral structure and made several sketches.
Though the galaxy can be detected with binoculars or a small telescope, to observe the spiral structure in a telescope without a camera requires a fairly large instrument, very dark skies, and a low-power eyepiece.
Structure and composition
M101 is a large galaxy, with a diameter of 170,000 light-years. By comparison, the Milky Way has a diameter of 87,400 light-years. It has around a trillion stars. It has a disk mass on the order of 100 billion solar masses, along with a small central bulge of about 3 billion solar masses. Its characteristics can be compared to those of Andromeda Galaxy.
M101 has a high population of H II regions, many of which are very large and bright. H II regions usually accompany the enormous clouds of high density molecular hydrogen gas contracting under their own gravitational force where stars form. H II regions are ionized by large numbers of extremely bright and hot young stars; those in M101 are capable of creating hot superbubbles. In a 1990 study, 1,264 H II regions were cataloged in the galaxy. Three are prominent enough to receive New General Catalogue numbers—NGC 5461, NGC 5462, and NGC 5471.
M101 is asymmetrical due to the tidal forces from interactions with its companion galaxies. These gravitational interactions compress interstellar hydrogen gas, which then triggers strong star formation activity in M101's spiral arms that can be detected in ultraviolet images.
In 2001, the X-ray source P98, located in M101, was identified as an ultra-luminous X-ray source—a source more powerful than any single star but less powerful than a whole galaxy—using the Chandra X-ray Observatory. It received the designation M101 ULX-1. In 2005, Hubble and XMM-Newton observations showed the presence of an optical counterpart, strongly indicating that M101 ULX-1 is an X-ray binary. Further observations showed that the system deviated from expected models—the black hole is just 20 to 30 solar masses, and consumes material (including captured stellar wind) at a higher rate than theory suggests.
It is estimated that M101 has about 150 globular clusters, the same as the number of the Milky Way's globular clusters.
Companion galaxies
M101 has six prominent companion galaxies: NGC 5204, NGC 5474, NGC 5477, NGC 5585, UGC 8837 and UGC 9405. As stated above, the gravitational interaction between it and its satellites may have spawned its grand design pattern. The galaxy has probably distorted the second-listed companion. The list comprises most or all of the M101 Group.
Supernovae and luminous red nova
Six internal supernovae have been recorded:
SN 1909A was discovered by Max Wolf in January 1909 and reached magnitude 12.1.
SN 1951H was discovered by Milton Humason on 1 September 1951 and reached magnitude 17.5.
SN 1970G (typeII, mag. 11.5) was discovered by Miklós Lovas on 30 July 1970.
On August 24, 2011, a Type Ia supernova, SN 2011fe, initially designated PTF 11kly, was discovered in M101. It had visual magnitude 17.2 at discovery and reached 9.9 at its peak.
On February 10, 2015, a luminous red nova, known as M101 OT2015-1 was discovered in the Pinwheel Galaxy.
On May 19, 2023, SN 2023ixf was discovered in M101, and immediately classified as a Type II supernova.
| Physical sciences | Notable galaxies | null |
319126 | https://en.wikipedia.org/wiki/Lenticular%20galaxy | Lenticular galaxy | A lenticular galaxy (denoted S0) is a type of galaxy intermediate between an elliptical (denoted E) and a spiral galaxy in galaxy morphological classification schemes. It contains a large-scale disc but does not have large-scale spiral arms. Lenticular galaxies are disc galaxies that have used up or lost most of their interstellar matter and therefore have very little ongoing star formation. They may, however, retain significant dust in their disks. As a result, they consist mainly of aging stars (like elliptical galaxies). Despite the morphological differences, lenticular and elliptical galaxies share common properties like spectral features and scaling relations. Both can be considered early-type galaxies that are passively evolving, at least in the local part of the Universe. Connecting the E galaxies with the S0 galaxies are the ES galaxies with intermediate-scale discs.
Morphology and structure
Classification
Lenticular galaxies are unique in that they have a visible disk component as well as a prominent bulge component. They have much higher bulge-to-disk ratios than typical spirals and do not have the canonical spiral arm structure of late-type galaxies, yet may exhibit a central bar. This bulge dominance can be seen in the axis ratio (i.e. the ratio between the observed minor and major axial of a disk galaxy) distribution of a lenticular galaxy sample. The distribution for lenticular galaxies rises steadily in the range 0.25 to 0.85 whereas the distribution for spirals is essentially flat in that same range. Larger axial ratios can be explained by observing face-on disk galaxies or by having a sample of spheroidal (bulge-dominated) galaxies. Imagine looking at two disk galaxies edge-on, one with a bulge and one without a bulge. The galaxy with a prominent bulge will have a larger edge-on axial ratio compared to the galaxy without a bulge based on the definition of axial ratio. Thus a sample of disk galaxies with prominent spheroidal components will have more galaxies at larger axial ratios. The fact that the lenticular galaxy distribution rises with increasing observed axial ratio implies that lenticulars are dominated by a central bulge component.
Lenticular galaxies are often considered to be a poorly understood transition state between spiral and elliptical galaxies, which results in their intermediate placement on the Hubble sequence. This results from lenticulars having both prominent disk and bulge components. The disk component is usually featureless, which precludes a classification system similar to spiral galaxies. As the bulge component is usually spherical, elliptical galaxy classifications are also unsuitable. Lenticular galaxies are thus divided into subclasses based upon either the amount of dust present or the prominence of a central bar. The classes of lenticular galaxies with no bar are S01, S02, and S03 where the subscripted numbers indicate the amount of dust absorption in the disk component; the corresponding classes for lenticulars with a central bar are SB01, SB02, and SB03.
Sérsic decomposition
The surface brightness profiles of lenticular galaxies are well described by the sum of a Sérsic model for the spheroidal component plus an exponentially declining model (Sérsic index of n ≈ 1) for the disk, and often a third component for the bar. Sometimes there is an observed truncation in the surface brightness profiles of lenticular galaxies at ~ 4 disk scalelengths. These features are consistent with the general structure of spiral galaxies. However, the bulge component of lenticulars is more closely related to elliptical galaxies in terms of morphological classification. This spheroidal region, which dominates the inner structure of lenticular galaxies, has a steeper surface brightness profile (Sérsic index typically ranging from n = 1 to 4) than the disk component. Lenticular galaxy samples are distinguishable from the diskless (excluding small nuclear disks) elliptical galaxy population through analysis of their surface brightness profiles.
Bars
Like spiral galaxies, lenticular galaxies can possess a central bar structure. While the classification system for normal lenticulars depends on dust content, barred lenticular galaxies are classified by the prominence of the central bar. SB01 galaxies have the least defined bar structure and are only classified as having slightly enhanced surface brightness along opposite sides of the central bulge. The prominence of the bar increases with index number, thus SB03 galaxies, like the NGC 1460 have very well defined bars that can extend through the transition region between the bulge and disk. NGC 1460 is actually the galaxy with one of the largest bars seen among lenticular galaxies. Unfortunately, the properties of bars in lenticular galaxies have not been researched in great detail. Understanding these properties, as well as understanding the formation mechanism for bars, would help clarify the formation or evolution history of lenticular galaxies.
Box-shaped bulges
NGC 1375 and NGC 1175 are examples of lenticular galaxies that have so-called box-shaped bulges. They are classified as SB0 pec. Box-shaped bulges are seen in edge-on galaxies, mostly spiral, but rarely lenticular.
Content
In many respects the composition of lenticular galaxies is like that of ellipticals. For example, they both consist of predominately older, hence redder, stars. All of their stars are thought to be older than about a billion years, in agreement with their offset from the Tully–Fisher relation (see below). In addition to these general stellar attributes, globular clusters are found more frequently in lenticular galaxies than in spiral galaxies of similar mass and luminosity. They also have little to no molecular gas (hence the lack of star formation) and no significant hydrogen α or 21-cm emission. Finally, unlike ellipticals, they may still possess significant dust.
Kinematics
Measurement difficulties and techniques
Lenticular galaxies share kinematic properties with both spiral and elliptical galaxies. This is due to the significant bulge and disk nature of lenticulars. The bulge component is similar to elliptical galaxies in that it is pressure supported by a central velocity dispersion. This situation is analogous to a balloon, where the motions of the air particles (stars in a bulge's case) are dominated by random motions. However, the kinematics of lenticular galaxies are dominated by the rotationally supported disk. Rotation support implies the average circular motion of stars in the disk is responsible for the stability of the galaxy. Thus, kinematics are often used to distinguish lenticular galaxies from elliptical galaxies. Determining the distinction between elliptical galaxies and lenticular galaxies often relies on the measurements of velocity dispersion (σ), rotational velocity (v), and ellipticity (ε). In order to differentiate between lenticulars and ellipticals, one typically looks at the v/σ ratio for a fixed ε. For example, a rough criterion for distinguishing between lenticular and elliptical galaxies is that elliptical galaxies have v/σ < 0.5 for ε = 0.3. The motivation behind this criterion is that lenticular galaxies do have prominent bulge and disk components whereas elliptical galaxies have no disk structure. Thus, lenticulars have much larger v/σ ratios than ellipticals due to their non-negligible rotational velocities (due to the disk component) in addition to not having as prominent of a bulge component compared to elliptical galaxies. However, this approach using a single ratio for each galaxy is problematic due to the dependence of the v/σ ratio on the radius out to which it is measured in some early-type galaxies. For example, the ES galaxies that bridge the E and S0 galaxies, with their intermediate-scale disks, have a high v/σ ratio at intermediate radii that then drops to a low ratio at large radii.
The kinematics of disk galaxies are usually determined by Hα or 21-cm emission lines, which are typically not present in lenticular galaxies due to their general lack of cool gas. Thus kinematic information and rough mass estimates for lenticular galaxies often comes from stellar absorption lines, which are less reliable than emission line measurements. There is also a considerable amount of difficulty in deriving accurate rotational velocities for lenticular galaxies. This is a combined effect from lenticulars having difficult inclination measurements, projection effects in the bulge-disk interface region, and the random motions of stars affecting the true rotational velocities. These effects make kinematic measurements of lenticular galaxies considerably more difficult compared to normal disk galaxies.
Offset Tully–Fisher relation
The kinematic connection between spiral and lenticular galaxies is most clear when analyzing the Tully–Fisher relation for spiral and lenticular samples. If lenticular galaxies are an evolved stage of spiral galaxies then they should have a similar Tully–Fisher relation with spirals, but with an offset in the luminosity / absolute magnitude axis. This would result from brighter, redder stars dominating the stellar populations of lenticulars. An example of this effect can be seen in the adjacent plot. One can clearly see that the best-fit lines for the spiral galaxy data and the lenticular galaxy have the same slope (and thus follow the same Tully–Fisher relation), but are offset by ΔI ≈ 1.5. This implies that lenticular galaxies were once spiral galaxies but are now dominated by old, red stars.
Formation theories
The morphology and kinematics of lenticular galaxies each, to a degree, suggest a mode of galaxy formation. Their disk-like, possibly dusty, appearance suggests they come from faded spiral galaxies, whose arm features disappeared. However, some lenticular galaxies are more luminous than spiral galaxies, which suggests that they are not merely the faded remnants of spiral galaxies. Lenticular galaxies might result from a galaxy merger, which increase the total stellar mass and might give the newly merged galaxy a disk-like, arm-less appearance. Alternatively, it has been proposed that they grew their disks via (gas and minor merger) accretion events.
It had previously been suggested that the evolution of luminous lenticular galaxies may be closely linked to that of elliptical galaxies, whereas fainter lenticulars might be more closely associated with ram-pressure stripped spiral galaxies, although this latter galaxy harassment scenario has since been queried due to the existence of extremely isolated, low-luminosity lenticular galaxies such as LEDA 2108986.
Faded spirals
The absence of gas, presence of dust, lack of recent star formation, and rotational support are all attributes one might expect of a spiral galaxy which had used up all of its gas in the formation of stars. This possibility is further enhanced by the existence of gas poor, or "anemic", spiral galaxies. If the spiral pattern then dissipated the resulting galaxy would be similar to many lenticulars. Moore et al. also document that tidal harassment – the gravitational effects from other, near-by galaxies – could aid this process in dense regions. The clearest support for this theory, however, is their adherence to slightly shifted version of Tully–Fisher relation, discussed above.
A 2012 paper that suggests a new classification system, first proposed by the Canadian astronomer Sidney van den Bergh, for lenticular and dwarf spheroidal galaxies (S0a-S0b-S0c-dSph) that parallels the Hubble sequence for spirals and irregulars (Sa-Sb-Sc-Im) reinforces this idea showing how the spiral–irregular sequence is very similar to this new one for lenticulars and dwarf ellipticals.
Mergers
The analyses of Burstein and Sandage showed that lenticular galaxies typically have surface brightness much greater than other spiral classes. It is also thought that lenticular galaxies exhibit a larger bulge-to-disk ratio than spiral galaxies and this may be inconsistent with simple fading from a spiral. If S0s were formed by mergers of other spirals these observations would be fitting and it would also account for the increased frequency of globular clusters. It should be mentioned, however, that advanced models of the central bulge which include both a general Sersic profile and bar indicate a smaller bulge, and thus a lessened inconsistency. Mergers are also unable to account for the offset from the Tully–Fisher relation without assuming that the merged galaxies were quite different from those we see today.
Disk growth via accretion
The creation of disks in, at least some, lenticular galaxies via the accretion of gas, and small galaxies, around a pre-existing spheroidal structure was first suggested as an explanation to match the high-redshift compact massive spheroidal-shaped galaxies with the equally compact massive bulges seen in nearby massive lenticular galaxies. In a "downsizing" scenario, bigger lenticular galaxies may have been built first – in a younger universe when more gas was available – and the lower-mass galaxies may have been slower to attract their disk-building material, as in the case of the isolated early-type galaxy LEDA 2108986. Within galaxy clusters, ram-pressure stripping removes gas and prevents the accretion of new gas that might be capable of furthering the development of the disk.
Examples
Cartwheel Galaxy, lenticular galaxy about 500 million light-years away in the constellation Sculptor
NGC 2787, a barred lenticular galaxy
NGC 3115
NGC 3632
NGC 4608, a barred lenticular galaxy about 56 million light years away in Virgo
NGC 5866
NGC 1533 is a prototypical lenticular galaxy in the constellation Dorado
Gallery
| Physical sciences | Galaxy morphological classification | null |
319141 | https://en.wikipedia.org/wiki/Eagle%20Nebula | Eagle Nebula | The Eagle Nebula (catalogued as Messier 16 or M16, and as NGC 6611, and also known as the Star Queen Nebula) is a young open cluster of stars in the constellation Serpens, discovered by Jean-Philippe de Cheseaux in 1745–46. Both the "Eagle" and the "Star Queen" refer to visual impressions of the dark silhouette near the center of the nebula, an area made famous as the "Pillars of Creation" imaged by the Hubble Space Telescope. The nebula contains several active star-forming gas and dust regions, including the aforementioned Pillars of Creation. The Eagle Nebula lies in the Sagittarius Arm of the Milky Way.
Characteristics
The Eagle Nebula is part of a diffuse emission nebula, or H II region, which is catalogued as IC 4703. This region of active current star formation is about 5700 light-years distant. A spire of gas that can be seen coming off the nebula in the northeastern part is approximately 9.5 light-years or about 90 trillion kilometers long.
The cluster associated with the nebula has approximately 8100 stars, which are mostly concentrated in a gap in the molecular cloud to the north-west of the Pillars.
The brightest star (HD 168076) has an apparent magnitude of +8.24, easily visible with good binoculars. It is actually a binary star formed of an O3.5V star plus an O7.5V companion. This star has a mass of roughly 80 solar masses, and a luminosity up to 1 million times that of the Sun.
The cluster's age has been estimated to be 1–2 million years.
The descriptive names reflect impressions of the shape of the central pillar rising from the southeast into the central luminous area. The name "Star Queen Nebula" was introduced by Robert Burnham, Jr., reflecting his characterization of the central pillar as the Star Queen shown in silhouette.
"Pillars of Creation" region
Images produced by Jeff Hester and Paul Scowen using the Hubble Space Telescope in 1995 greatly improved scientific understanding of processes inside the nebula. One of these became famous as the "Pillars of Creation", depicting a large region of star formation. Its small dark pockets are believed to be protostars (Bok globules). The pillar structure resembles that of a much larger instance in the Soul Nebula of Cassiopeia, imaged with the Spitzer Space Telescope in 2005 equally characterized as "Pillars of Star Creation". or "Pillars of Star Formation". These columns – which resemble stalagmites protruding from the floor of a cavern – are composed of interstellar hydrogen gas and dust, which act as incubators for new stars. Inside the columns and on their surface astronomers have found knots or globules of denser gas, called EGGs ("Evaporating Gaseous Globules"). Stars are being formed inside some of these.
X-ray images from the Chandra observatory compared with Hubble's "Pillars" image have shown that X-ray sources (from young stars) do not coincide with the pillars, but rather randomly dot the nebula. Any protostars in the pillars' EGGs are not yet hot enough to emit X-rays.
Evidence from the Spitzer Space Telescope originally suggested that the pillars in M16 may be threatened by a "past supernova". Hot gas observed by Spitzer in 2007 suggested they were already – likely – being disturbed by a supernova that exploded 8,000 to 9,000 years ago. Due to the distance the main blast of light would have reached Earth for a brief time 1,000 to 2,000 years ago. A more slowly moving, theorized, shock wave would have taken a few thousand years to move through the nebula and would have blown away the delicate pillars. However, in 2014 the pillars were imaged a second time by Hubble, in both visible light and infrared light. The images being 20 years later provided a new, detailed account of the rate of evaporation occurring within the pillars. No supernova is evidenced within them, and it is estimated in some form they still exist – and will appear for at least 100,000 more years.
Gallery
| Physical sciences | Notable nebulae | null |
319150 | https://en.wikipedia.org/wiki/Messier%204 | Messier 4 | Messier 4 or M4 (also known as NGC 6121 or the Spider Globular Cluster) is a globular cluster in the constellation of Scorpius. It was discovered by Philippe Loys de Chéseaux in 1745 and catalogued by Charles Messier in 1764. It was the first globular cluster in which individual stars were resolved.
Visibility
M4 is conspicuous in even the smallest of telescopes as a fuzzy ball of light. It appears about the same size as the Moon in the sky. It is one of the easiest globular clusters to find, being located only 1.3 degrees west of the bright star Antares, with both objects being visible in a wide-field telescope. Modestly sized telescopes will begin to resolve individual stars, of which the brightest in M4 are of apparent magnitude 10.8.
Characteristics
M4 is a rather loosely concentrated cluster of class IX and measures 75 light-years across. It features a characteristic "bar" structure across its core, visible to moderate sized telescopes. The structure consists of 11th-magnitude stars and is approximately 2.5' long and was first noted by William Herschel in 1783. At least 43 variable stars have been observed within M4.
M4 is approximately 6,000 light-years away, making it the closest globular cluster to the Solar System. It has an estimated age of 12.2 billion years.
In astronomy, the abundance of elements other than hydrogen and helium is called the metallicity, and it is usually denoted by the abundance ratio of iron to hydrogen as compared to the Sun. For this cluster, the measured abundance of iron is equal to:
This value is the logarithm of the ratio of iron to hydrogen relative to the same ratio in the Sun. Thus the cluster has an abundance of iron equal to 8.5% of the iron abundance in the Sun. This strongly suggests this cluster hosts two distinct stellar populations, differing by age. Thus the cluster probably saw two main cycles or phases of star formation.
The space velocity components are (U, V, W) = (, , ) km/s. This confirms an orbit around the Milky Way of a period of with eccentricity 0.80 ± 0.03: during periapsis it comes within from the galactic core, while at apoapsis it travels out to . The inclination is at (an angle of) from the galactic plane, thus it reaches as much as above the disk. When passing through the disk, this cluster does so at less than 5 kpc from the galactic nucleus. The cluster undergoes tidal shock during each passage, which can cause the repeated shedding of stars. Thus the cluster may have been much more massive.
Notable stars
Photographs by the Hubble Space Telescope in 1995 found white dwarf stars in M4 that are among the oldest known stars in our galaxy; aged 13 billion years. One has been found to be a binary star with a pulsar companion, PSR B1620−26 and a planet orbiting it with a mass of 2.5 times that of Jupiter (). One star in Messier 4 was also found to have much more of the rare light element lithium than expected.
CX-1 Is located in M4. It is known as a possible millisecond pulsar/neutron star binary. It orbits in 6.31 hours.
Spinthariscope analogy
The view of Messier 4 through a good telescope was likened by Robert Burnham Jr. to that of hyperkinetic luminous alpha particles seen in a spinthariscope.
Central black hole
In 2023, an analysis of Hubble Space Telescope and European Space Agency's Gaia spacecraft data from Messier 4 revealed an excess mass of roughly 800 solar masses in the center of this cluster, which appears to not be extended. This could thus be considered as kinematic evidence for an intermediate-mass black hole (even if an unusually compact cluster of compact objects like white dwarfs, neutron stars or stellar-mass black holes cannot be completely discounted).
| Physical sciences | Notable star clusters | Astronomy |
319153 | https://en.wikipedia.org/wiki/Messier%2083 | Messier 83 | Messier 83 or M83, also known as the Southern Pinwheel Galaxy and NGC 5236, is a barred spiral galaxy approximately 15 million light-years away in the constellation borders of Hydra and Centaurus. Nicolas-Louis de Lacaille discovered M83 on 17 February 1752 at the Cape of Good Hope. Charles Messier added it to his catalogue of nebulous objects (now known as the Messier Catalogue) in March 1781.
It is one of the closest and brightest barred spiral galaxies in the sky, and is visible with binoculars. It has an isophotal diameter at about . Its nickname of the Southern Pinwheel derives from its resemblance to the Pinwheel Galaxy (M101).
Characteristics
M83 is a massive, grand design spiral galaxy. Its morphological classification in the De Vaucouleurs system is SAB(s)c, where the 'SAB' denotes a weak-barred spiral, '(s)' indicates a pure spiral structure with no ring, and 'c' means the spiral arms are loosely wound. The peculiar dwarf galaxy NGC 5253 lies near M83, and the two likely interacted within the last billion years resulting in starburst activity in their central regions.
The star formation rate in M83 is higher along the leading edge of the spiral arms, as predicted by density wave theory. NASA's Galaxy Evolution Explorer project on 16 April 2008 reported finding large numbers of new stars in the outer reaches of the galaxy— from the center. It had been thought that these areas lacked the materials necessary for star formation.
Supernovae
Six supernovae have been observed in M83:
SN 1923A (type unknown, mag. 14) was discovered by Carl Otto Lampland on 5 May 1923.
SN 1945B (type unknown, mag. 14.2) was discovered by William Liller on 13 July 1945.
SN 1950B (type unknown, mag. 14.5) was discovered by Guillermo Haro on 15 March 1950.
SN 1957D (type unknown, mag. 15) was discovered by H. S. Gates on 28 December 1957.
SN 1968L (type II-P, mag. 11.9) was discovered by J. C. Bennett on 17 July 1968.
SN 1983N (type Ia, mag. 11.9) was discovered by Robert Evans from Australia on July 3, 1983. On July 6, it was observed with the Very Large Array and became the first type I supernova to have a radio emission detected. The supernova reached peak optical brightness on July 17, achieving an apparent visual magnitude of 11.54. Although identified as type I, the spectrum was considered peculiar. A year after the explosion, about of iron was discovered in the ejecta. This was the first time that such a large amount of iron was unambiguously detected from a supernova explosion. SN 1983N became the modern prototype of a hydrogen deficient type Ib supernova, with the progenitor being inferred as a Wolf–Rayet star.
Environment
M83 is at the center of one of two subgroups within the Centaurus A/M83 Group, a nearby galaxy group. Centaurus A is at the center of the other subgroup. These are sometimes identified as one group, and sometimes as two. However, the galaxies around Centaurus A and the galaxies around M83 are physically close to each other, and both subgroups appear not to be moving relative to each other.
| Physical sciences | Notable galaxies | Astronomy |
319168 | https://en.wikipedia.org/wiki/Shield%20volcano | Shield volcano | A shield volcano is a type of volcano named for its low profile, resembling a shield lying on the ground. It is formed by the eruption of highly fluid (low viscosity) lava, which travels farther and forms thinner flows than the more viscous lava erupted from a stratovolcano. Repeated eruptions result in the steady accumulation of broad sheets of lava, building up the shield volcano's distinctive form.
Shield volcanoes are found wherever fluid, low-silica lava reaches the surface of a rocky planet. However, they are most characteristic of ocean island volcanism associated with hot spots or with continental rift volcanism. They include the largest active volcanoes on Earth, such as Mauna Loa. Giant shield volcanoes are found on other planets of the Solar System, including Olympus Mons on Mars and Sapas Mons on Venus.
Etymology
The term 'shield volcano' is taken from the German term Schildvulkan, coined by the Austrian geologist Eduard Suess in 1888 and which had been calqued into English by 1910.
Geology
Structure
Shield volcanoes are distinguished from the three other major volcanic types—stratovolcanoes, lava domes, and cinder cones—by their structural form, a consequence of their particular magmatic composition. Of these four forms, shield volcanoes erupt the least viscous lavas. Whereas stratovolcanoes and lava domes are the product of highly viscous flows, and cinder cones are constructed of explosively eruptive tephra, shield volcanoes are the product of gentle effusive eruptions of highly fluid lavas that produce, over time, a broad, gently sloped eponymous "shield". Although the term is generally applied to basaltic shields, it has also at times been applied to rarer scutiform volcanoes of differing magmatic composition—principally pyroclastic shields, formed by the accumulation of fragmentary material from particularly powerful explosive eruptions, and rarer felsic lava shields formed by unusually fluid felsic magmas. Examples of pyroclastic shields include Billy Mitchell volcano in Papua New Guinea and the Purico complex in Chile; an example of a felsic shield is the Ilgachuz Range in British Columbia, Canada. Shield volcanoes are similar in origin to vast lava plateaus and flood basalts present in various parts of the world. These are eruptive features which occur along linear fissure vents and are distinguished from shield volcanoes by the lack of an identifiable primary eruptive center.
Active shield volcanoes experience near-continuous eruptive activity over extremely long periods of time, resulting in the gradual build-up of edifices that can reach extremely large dimensions. With the exclusion of flood basalts, mature shields are the largest volcanic features on Earth. The summit of the largest subaerial volcano in the world, Mauna Loa, lies above sea level, and the volcano, over wide at its base, is estimated to contain about of basalt. The mass of the volcano is so great that it has slumped the crust beneath it a further . Accounting for this subsidence and for the height of the volcano above the sea floor, the "true" height of Mauna Loa from the start of its eruptive history is about . Mount Everest, by comparison, is in height. In 2013, a team led by the University of Houston's William Sager announced the discovery of Tamu Massif, an enormous extinct submarine volcano, approximately in area, which dwarfs all previously known volcanoes on Earth. However, the extents of the volcano have not been confirmed. Although Tamu Massif was initially believed to be a shield volcano, Sanger and his colleagues acknowledged in 2019 that Tamu Massif is not a shield volcano.
Shield volcanoes feature a gentle (usually 2° to 3°) slope that gradually steepens with elevation (reaching approximately 10°) before flattening near the summit, forming an overall upwardly convex shape. These slope characteristics have a correlation with age of the forming lava, with in the case of the Hawaiian chain, steepness increasing with age, as later lavas tend to be more alkali so are more viscous, with thicker flows, that travel less distance from the summit vents. In height they are typically about one twentieth their width. Although the general form of a "typical" shield volcano varies little worldwide, there are regional differences in their size and morphological characteristics. Typical shield volcanoes found in California and Oregon measure in diameter and in height, while shield volcanoes in the central Mexican Michoacán–Guanajuato volcanic field average in height and in width, with an average slope angle of 9.4° and an average volume of .
Rift zones are a prevalent feature on shield volcanoes that is rare on other volcanic types. The large, decentralized shape of Hawaiian volcanoes as compared to their smaller, symmetrical Icelandic cousins can be attributed to rift eruptions. Fissure venting is common in Hawaii; most Hawaiian eruptions begin with a so-called "wall of fire" along a major fissure line before centralizing to a small number of points. This accounts for their asymmetrical shape, whereas Icelandic volcanoes follow a pattern of central eruptions dominated by summit calderas, causing the lava to be more evenly distributed or symmetrical.
Eruptive characteristics
Most of what is currently known about shield volcanic eruptive character has been gleaned from studies done on the volcanoes of Hawaii Island, by far the most intensively studied of all shields because of their scientific accessibility; the island lends its name to the slow-moving, effusive eruptions typical of shield volcanism, known as Hawaiian eruptions. These eruptions, the least explosive of volcanic events, are characterized by the effusive emission of highly fluid basaltic lavas with low gaseous content. These lavas travel a far greater distance than those of other eruptive types before solidifying, forming extremely wide but relatively thin magmatic sheets often less than thick. Low volumes of such lavas layered over long periods of time are what slowly constructs the characteristically low, broad profile of a mature shield volcano.
Also unlike other eruptive types, Hawaiian eruptions often occur at decentralized fissure vents, beginning with large "curtains of fire" that quickly die down and concentrate at specific locations on the volcano's rift zones. Central-vent eruptions, meanwhile, often take the form of large lava fountains (both continuous and sporadic), which can reach heights of hundreds of meters or more. The particles from lava fountains usually cool in the air before hitting the ground, resulting in the accumulation of cindery scoria fragments; however, when the air is especially thick with pyroclasts, they cannot cool off fast enough because of the surrounding heat, and hit the ground still hot, accumulating into spatter cones. If eruptive rates are high enough, they may even form splatter-fed lava flows. Hawaiian eruptions are often extremely long-lived; Puʻu ʻŌʻō, a cinder cone of Kīlauea, erupted continuously from January 3, 1983, until April 2018.
Flows from Hawaiian eruptions can be divided into two types by their structural characteristics: pāhoehoe lava which is relatively smooth and flows with a ropey texture, and ʻaʻā flows which are denser, more viscous (and thus slower moving) and blockier. These lava flows can be anywhere between thick. Aā lava flows move through pressure— the partially solidified front of the flow steepens because of the mass of flowing lava behind it until it breaks off, after which the general mass behind it moves forward. Though the top of the flow quickly cools down, the molten underbelly of the flow is buffered by the solidifying rock above it, and by this mechanism, aā flows can sustain movement for long periods of time. Pāhoehoe flows, in contrast, move in more conventional sheets, or by the advancement of lava "toes" in snaking lava columns. Increasing viscosity on the part of the lava or shear stress on the part of local topography can morph a pāhoehoe flow into an ʻaʻā one, but the reverse never occurs.
Although most shield volcanoes are by volume almost entirely Hawaiian and basaltic in origin, they are rarely exclusively so. Some volcanoes, such as Mount Wrangell in Alaska and Cofre de Perote in Mexico, exhibit large enough swings in their historical magmatic eruptive characteristics to cast strict categorical assignment in doubt; one geological study of de Perote went so far as to suggest the term "compound shield-like volcano" instead. Most mature shield volcanoes have multiple cinder cones on their flanks, the results of tephra ejections common during incessant activity and markers of currently and formerly active sites on the volcano. An example of these parasitic cones is at Puʻu ʻŌʻō on Kīlauea—continuous activity ongoing since 1983 has built up a tall cone at the site of one of the longest-lasting rift eruptions in known history.
The Hawaiian shield volcanoes are not located near any plate boundaries; the volcanic activity of this island chain is distributed by the movement of the oceanic plate over an upwelling of magma known as a hotspot. Over millions of years, the tectonic movement that moves continents also creates long volcanic trails across the seafloor. The Hawaiian and Galápagos shields, and other hotspot shields like them, are constructed of oceanic island basalt. Their lavas are characterized by high levels of sodium, potassium, and aluminium.
Features common in shield volcanism include lava tubes. Lava tubes are cave-like volcanic straights formed by the hardening of overlaying lava. These structures help further the propagation of lava, as the walls of the tube insulate the lava within. Lava tubes can account for a large portion of shield volcano activity; for example, an estimated 58% of the lava forming Kīlauea comes from lava tubes.
In some shield volcano eruptions, basaltic lava pours out of a long fissure instead of a central vent, and shrouds the countryside with a long band of volcanic material in the form of a broad plateau. Plateaus of this type exist in Iceland, Washington, Oregon, and Idaho; the most prominent ones are situated along the Snake River in Idaho and the Columbia River in Washington and Oregon, where they have been measured to be over in thickness.
Calderas are a common feature on shield volcanoes. They are formed and reformed over the volcano's lifespan. Long eruptive periods form cinder cones, which then collapse over time to form calderas. The calderas are often filled up by progressive eruptions, or formed elsewhere, and this cycle of collapse and regeneration takes place throughout the volcano's lifespan.
Interactions between water and lava at shield volcanoes can cause some eruptions to become hydrovolcanic. These explosive eruptions are drastically different from the usual shield volcanic activity and are especially prevalent at the waterbound volcanoes of the Hawaiian Isles.
Distribution
Shield volcanoes are found worldwide. They can form over hotspots (points where magma from below the surface wells up), such as the Hawaiian–Emperor seamount chain and the Galápagos Islands, or over more conventional rift zones, such as the Icelandic shields and the shield volcanoes of East Africa. Although shield volcanoes are not usually associated with subduction, they can occur over subduction zones. Many examples are found in California and Oregon, including Prospect Peak in Lassen Volcanic National Park, as well as Pelican Butte and Belknap Crater in Oregon. Many shield volcanoes are found in ocean basins, such as Kīlauea in Hawaii, although they can be found inland as well—East Africa being one example of this.
Hawaiian–Emperor seamount chain
The largest and most prominent shield volcano chain in the world is the Hawaiian–Emperor seamount chain, a chain of hotspot volcanoes in the Pacific Ocean. The volcanoes follow a distinct evolutionary pattern of growth and death. The chain contains at least 43 major volcanoes, and Meiji Seamount at its terminus near the Kuril–Kamchatka Trench is 85 million years old.
The youngest part of the chain is Hawaii, where the volcanoes are characterized by frequent rift eruptions, their large size (thousands of km3 in volume), and their rough, decentralized shape. Rift zones are a prominent feature on these volcanoes and account for their seemingly random volcanic structure. They are fueled by the movement of the Pacific Plate over the Hawaii hotspot and form a long chain of volcanoes, atolls, and seamounts long with a total volume of over .
The chain includes Mauna Loa, a shield volcano which stands above sea level and reaches a further below the waterline and into the crust, approximately of rock. Kīlauea, another Hawaiian shield volcano, is one of the most active volcanoes on Earth, with its most recent eruption occurring in 2021.
Galápagos Islands
The Galápagos Islands are an isolated set of volcanoes, consisting of shield volcanoes and lava plateaus, about west of Ecuador. They are driven by the Galápagos hotspot, and are between approximately 4.2 million and 700,000 years of age. The largest island, Isabela, consists of six coalesced shield volcanoes, each delineated by a large summit caldera. Española, the oldest island, and Fernandina, the youngest, are also shield volcanoes, as are most of the other islands in the chain. The Galápagos Islands are perched on a large lava plateau known as the Galápagos Platform. This platform creates a shallow water depth of at the base of the islands, which stretch over a diameter. Since Charles Darwin's visit to the islands in 1835 during the second voyage of HMS Beagle, there have been over 60 recorded eruptions in the islands, from six different shield volcanoes. Of the 21 emergent volcanoes, 13 are considered active.
Cerro Azul is a shield volcano on the southwestern part of Isabela Island and is one of the most active in the Galapagos, with the last eruption between May and June 2008. The Geophysics Institute at the National Polytechnic School in Quito houses an international team of seismologists and volcanologists whose responsibility is to monitor Ecuador's numerous active volcanoes in the Andean Volcanic Belt and the Galapagos Islands. La Cumbre is an active shield volcano on Fernandina Island that has been erupting since April 11, 2009.
The Galápagos islands are geologically young for such a big chain, and the pattern of their rift zones follows one of two trends, one north-northwest, and one east–west. The composition of the lavas of the Galápagos shields are strikingly similar to those of the Hawaiian volcanoes. Curiously, they do not form the same volcanic "line" associated with most hotspots. They are not alone in this regard; the Cobb–Eickelberg Seamount chain in the North Pacific is another example of such a delineated chain. In addition, there is no clear pattern of age between the volcanoes, suggesting a complicated, irregular pattern of creation. How the islands were formed remains a geological mystery, although several theories have been proposed.
Iceland
Located over the Mid-Atlantic Ridge, a divergent tectonic plate boundary in the middle of the Atlantic Ocean, Iceland is the site of about 130 volcanoes of various types. Icelandic shield volcanoes are generally of Holocene age, between 5,000 and 10,000 years old. The volcanoes are also very narrow in distribution, occurring in two bands in the West and North Volcanic Zones. Like Hawaiian volcanoes, their formation initially begins with several eruptive centers before centralizing and concentrating at a single point. The main shield then forms, burying the smaller ones formed by the early eruptions with its lava.
Icelandic shields are mostly small (~), symmetrical (although this can be affected by surface topography), and characterized by eruptions from summit calderas. They are composed of either tholeiitic olivine or picritic basalt. The tholeiitic shields tend to be wider and shallower than the picritic shields. They do not follow the pattern of caldera growth and destruction that other shield volcanoes do; caldera may form, but they generally do not disappear.
Turkey
Bingöl Mountains are one of the shield volcanoes in Turkey.
East Africa
In East Africa, volcanic activity is generated by the development of the East African Rift and from nearby hotspots. Some volcanoes interact with both. Shield volcanoes are found near the rift and off the coast of Africa, although stratovolcanoes are more common. Although sparsely studied, the fact that all of its volcanoes are of Holocene age reflects how young the volcanic center is. One interesting characteristic of East African volcanism is a penchant for the formation of lava lakes; these semi-permanent lava bodies, extremely rare elsewhere, form in about 9% of African eruptions.
The most active shield volcano in Africa is Nyamuragira. Eruptions at the shield volcano are generally centered within the large summit caldera or on the numerous fissures and cinder cones on the volcano's flanks. Lava flows from the most recent century extend down the flanks more than from the summit, reaching as far as Lake Kivu. Erta Ale in Ethiopia is another active shield volcano and one of the few places in the world with a permanent lava lake, which has been active since at least 1967, and possibly since 1906. Other volcanic centers include Menengai, a massive shield caldera, and Mount Marsabit in Kenya.
Extraterrestrial shield volcanoes
Shield volcanoes are not limited to Earth; they have been found on Mars, Venus, and Jupiter's moon, Io.
The shield volcanoes of Mars are very similar to the shield volcanoes on Earth. On both planets, they have gently sloping flanks, collapse craters along their central structure, and are built of highly fluid lavas. Volcanic features on Mars were observed long before they were first studied in detail during the 1976–1979 Viking mission. The principal difference between the volcanoes of Mars and those on Earth is in terms of size; Martian volcanoes range in size up to high and in diameter, far larger than the high, wide Hawaiian shields. The highest of these, Olympus Mons, is the tallest known mountain on any planet in the solar system.
Venus has over 150 shield volcanoes which are much flatter, with a larger surface area than those found on Earth, some having a diameter of more than . Although the majority of these are long extinct it has been suggested, from observations by the Venus Express spacecraft, that many may still be active.
| Physical sciences | Volcanology | Earth science |
319342 | https://en.wikipedia.org/wiki/Microfilament | Microfilament | Microfilaments, also called actin filaments, are protein filaments in the cytoplasm of eukaryotic cells that form part of the cytoskeleton. They are primarily composed of polymers of actin, but are modified by and interact with numerous other proteins in the cell. Microfilaments are usually about 7 nm in diameter and made up of two strands of actin. Microfilament functions include cytokinesis, amoeboid movement, cell motility, changes in cell shape, endocytosis and exocytosis, cell contractility, and mechanical stability. Microfilaments are flexible and relatively strong, resisting buckling by multi-piconewton compressive forces and filament fracture by nanonewton tensile forces. In inducing cell motility, one end of the actin filament elongates while the other end contracts, presumably by myosin II molecular motors. Additionally, they function as part of actomyosin-driven contractile molecular motors, wherein the thin filaments serve as tensile platforms for myosin's ATP-dependent pulling action in muscle contraction and pseudopod advancement. Microfilaments have a tough, flexible framework which helps the cell in movement.
Actin was first discovered in rabbit skeletal muscle in the mid 1940s by F.B. Straub. Almost 20 years later, H.E. Huxley demonstrated that actin is essential for muscle constriction. The mechanism in which actin creates long filaments was first described in the mid 1980s. Later studies showed that actin has an important role in cell shape, motility, and cytokinesis.
Organization
Actin filaments are assembled in two general types of structures: bundles and networks. Bundles can be composed of polar filament arrays, in which all barbed ends point to the same end of the bundle, or non-polar arrays, where the barbed ends point towards both ends. A class of actin-binding proteins, called cross-linking proteins, dictate the formation of these structures. Cross-linking proteins determine filament orientation and spacing in the bundles and networks. These structures are regulated by many other classes of actin-binding proteins, including motor proteins, branching proteins, severing proteins, polymerization promoters, and capping proteins.
In vitro self-assembly
Measuring approximately 6 nm in diameter, microfilaments are the thinnest fibers of the cytoskeleton. They are polymers of actin subunits (globular actin, or G-actin), which as part of the fiber are referred to as filamentous actin, or F-actin. Each microfilament is made up of two helical, interlaced strands of subunits. Much like microtubules, actin filaments are polarized. Electron micrographs have provided evidence of their fast-growing barbed-ends and their slow-growing pointed-end. This polarity has been determined by the pattern created by the binding of myosin S1 fragments: they themselves are subunits of the larger myosin II protein complex. The pointed end is commonly referred to as the minus (−) end and the barbed end is referred to as the plus (+) end.
In vitro actin polymerization, or nucleation, starts with the self-association of three G-actin monomers to form a trimer. ATP-bound actin then itself binds the barbed end, and the ATP is subsequently hydrolyzed. ATP hydrolysis occurs with a half time of about 2 seconds, while the half time for the dissociation of the inorganic phosphate is about 6 minutes. This autocatalyzed event reduces the binding strength between neighboring subunits, and thus generally destabilizes the filament. In vivo actin polymerization is catalyzed by a class of filament end-tracking molecular motors known as actoclampins. Recent evidence suggests that the rate of ATP hydrolysis and the rate of monomer incorporation are strongly coupled.
Subsequently, ADP-actin dissociates slowly from the pointed end, a process significantly accelerated by the actin-binding protein, cofilin. ADP bound cofilin severs ADP-rich regions nearest the (−)-ends. Upon release, the free actin monomer slowly dissociates from ADP, which in turn rapidly binds to the free ATP diffusing in the cytosol, thereby forming the ATP-actin monomeric units needed for further barbed-end filament elongation. This rapid turnover is important for the cell's movement. End-capping proteins such as CapZ prevent the addition or loss of monomers at the filament end where actin turnover is unfavorable, such as in the muscle apparatus.
Actin polymerization together with capping proteins were recently used to control the 3-dimensional growth of protein filament so as to perform 3D topologies useful in technology and the making of electrical interconnect. Electrical conductivity is obtained by metallisation of the protein 3D structure.
Mechanism of force generation
As a result of ATP hydrolysis, filaments elongate approximately 10 times faster at their barbed ends than their pointed ends. At steady-state, the polymerization rate at the barbed end matches the depolymerization rate at the pointed end, and microfilaments are said to be treadmilling. Treadmilling results in elongation in the barbed end and shortening in the pointed-end, so that the filament in total moves. Since both processes are energetically favorable, this means force is generated, the energy ultimately coming from ATP.
Actin in cells
Intracellular actin cytoskeletal assembly and disassembly are tightly regulated by cell signaling mechanisms. Many signal transduction systems use the actin cytoskeleton as a scaffold, holding them at or near the inner face of the peripheral membrane. This subcellular location allows immediate responsiveness to transmembrane receptor action and the resulting cascade of signal-processing enzymes.
Because actin monomers must be recycled to sustain high rates of actin-based motility during chemotaxis, cell signalling is believed to activate cofilin, the actin-filament depolymerizing protein which binds to ADP-rich actin subunits nearest the filament's pointed-end and promotes filament fragmentation, with concomitant depolymerization in order to liberate actin monomers. In most animal cells, monomeric actin is bound to profilin and thymosin beta-4, both of which preferentially bind with one-to-one stoichiometry to ATP-containing monomers. Although thymosin beta-4 is strictly a monomer-sequestering protein, the behavior of profilin is far more complex. Profilin enhances the ability of monomers to assemble by stimulating the exchange of actin-bound ADP for solution-phase ATP to yield actin-ATP and ADP. Profilin is transferred to the leading edge by virtue of its PIP2 binding site, and it employs its poly-L-proline binding site to dock onto end-tracking proteins. Once bound, profilin-actin-ATP is loaded into the monomer-insertion site of actoclampin motors.
Another important component in filament formation is the Arp2/3 complex, which binds to the side of an already existing filament (or "mother filament"), where it nucleates the formation of a new daughter filament at a 70-degree angle relative to the mother filament, effecting a fan-like branched filament network.
Specialized unique actin cytoskeletal structures are found adjacent to the plasma membrane. Four remarkable examples include red blood cells, human embryonic kidney cells, neurons, and sperm cells. In red blood cells, a spectrin-actin hexagonal lattice is formed by interconnected short actin filaments. In human embryonic kidney cells, the cortical actin forms a scale-free fractal structure. First found in neuronal axons, actin forms periodic rings that are stabilized by spectrin and adducin and this ring structure was then found by He et al 2016 to occur in almost every neuronal type and glial cells, across seemingly every animal taxon including Caenorhabditis elegans, Drosophila, Gallus gallus and Mus musculus. And in mammalian sperm, actin forms a helical structure in the midpiece, i.e., the first segment of the flagellum.
Associated proteins
In non-muscle cells, actin filaments are formed proximal to membrane surfaces. Their formation and turnover are regulated by many proteins, including:
Filament end-tracking protein (e.g., formins, VASP, N-WASP)
Filament-nucleator known as the Actin-Related Protein-2/3 (or Arp2/3) complex
Filament cross-linkers (e.g., α-actinin, fascin, and fimbrin)
Actin monomer-binding proteins profilin and thymosin β4
Filament barbed-end cappers such as Capping Protein and CapG, etc.
Filament-severing proteins like gelsolin.
Actin depolymerizing proteins such as ADF/cofilin.
The actin filament network in non-muscle cells is highly dynamic. The actin filament network is arranged with the barbed-end of each filament attached to the cell's peripheral membrane by means of clamped-filament elongation motors, the above-mentioned "actoclampins", formed from a filament barbed-end and a clamping protein (formins, VASP, Mena, WASP, and N-WASP). The primary substrate for these elongation motors is profilin-actin-ATP complex which is directly transferred to elongating filament ends. The pointed-end of each filament is oriented toward the cell's interior. In the case of lamellipodial growth, the Arp2/3 complex generates a branched network, and in filopodia a parallel array of filaments is formed.
Actin acts as a track for myosin motor motility
Myosin motors are intracellular ATP-dependent enzymes that bind to and move along actin filaments. Various classes of myosin motors have very different behaviors, including exerting tension in the cell and transporting cargo vesicles.
A proposed model – actoclampins track filament ends
One proposed model suggests the existence of actin filament barbed-end-tracking molecular motors termed "actoclampin". The proposed actoclampins generate the propulsive forces needed for actin-based motility of lamellipodia, filopodia, invadipodia, dendritic spines, intracellular vesicles, and motile processes in endocytosis, exocytosis, podosome formation, and phagocytosis. Actoclampin motors also propel such intracellular pathogens as Listeria monocytogenes, Shigella flexneri, Vaccinia and Rickettsia. When assembled under suitable conditions, these end-tracking molecular motors can also propel biomimetic particles.
The term actoclampin is derived from acto- to indicate the involvement of an actin filament, as in actomyosin, and clamp to indicate a clasping device used for strengthening flexible/moving objects and for securely fastening two or more components, followed by the suffix -in to indicate its protein origin. An actin filament end-tracking protein may thus be termed a clampin.
Dickinson and Purich recognized that prompt ATP hydrolysis could explain the forces achieved during actin-based motility. They proposed a simple mechanoenzymatic sequence known as the Lock, Load & Fire Model, in which an end-tracking protein remains tightly bound ("locked" or clamped) onto the end of one sub-filament of the double-stranded actin filament. After binding to Glycyl-Prolyl-Prolyl-Prolyl-Prolyl-Prolyl-registers on tracker proteins, Profilin-ATP-actin is delivered ("loaded") to the unclamped end of the other sub-filament, whereupon ATP within the already clamped terminal subunit of the other subfragment is hydrolyzed ("fired"), providing the energy needed to release that arm of the end-tracker, which then can bind another Profilin-ATP-actin to begin a new monomer-addition round.
Steps involved
The following steps describe one force-generating cycle of an actoclampin molecular motor:
The polymerization cofactor profilin and the ATP·actin combine to form a profilin-ATP-actin complex that then binds to the end-tracking unit
The cofactor and monomer are transferred to the barbed-end of an actin already clamped filament
The tracking unit and cofactor dissociate from the adjacent protofilament, in a step that can be facilitated by ATP hydrolysis energy to modulate the affinity of the cofactor and/or the tracking unit for the filament; and this mechanoenzymatic cycle is then repeated, starting this time on the other sub-filament growth site.
When operating with the benefit of ATP hydrolysis, AC motors generate per-filament forces of 8–9 pN, which is far greater than the per-filament limit of 1–2 pN for motors operating without ATP hydrolysis. The term actoclampin is generic and applies to all actin filament end-tracking molecular motors, irrespective of whether they are driven actively by an ATP-activated mechanism or passively.
Some actoclampins (e.g., those involving Ena/VASP proteins, WASP, and N-WASP) apparently require Arp2/3-mediated filament initiation to form the actin polymerization nucleus that is then "loaded" onto the end-tracker before processive motility can commence. To generate a new filament, Arp2/3 requires a "mother" filament, monomeric ATP-actin, and an activating domain from Listeria ActA or the VCA region of N-WASP. The Arp2/3 complex binds to the side of the mother filament, forming a Y-shaped branch having a 70-degree angle with respect to the longitudinal axis of the mother filament. Then upon activation by ActA or VCA, the Arp complex is believed to undergo a major conformational change, bringing its two actin-related protein subunits near enough to each other to generate a new filament gate. Whether ATP hydrolysis may be required for nucleation and/or Y-branch release is a matter under active investigation.
| Biology and health sciences | Cell parts | Biology |
319522 | https://en.wikipedia.org/wiki/Cattle%20egret | Cattle egret | The cattle egret (formerly genus Bubulcus) is a cosmopolitan clade of heron (family Ardeidae) in the genus Ardea found in the tropics, subtropics, and warm-temperate zones. According to the IOC bird list, it contains two species, the western cattle egret and the eastern cattle egret, although some authorities regard them as a single species. Despite the similarities in plumage to the egrets of the genus Egretta, it actually belongs to the genus Ardea. Originally native to parts of Asia, Africa, and Europe, it has undergone a rapid expansion in its distribution and successfully colonised much of the rest of the world in the last century.
They are white birds adorned with buff plumes in the breeding season. They nest in colonies, usually near bodies of water and often with other wading birds. The nest is a platform of sticks in trees or shrubs. Cattle egrets exploit drier and open habitats more than other heron species. Their feeding habitats include seasonally inundated grasslands, pastures, farmlands, wetlands, and rice paddies. They often accompany cattle or other large mammals, catching insect and small vertebrate prey disturbed by these animals. Some populations are migratory and others show postbreeding dispersal.
The adult cattle egret has few predators, but birds or mammals may raid its nests, and chicks may be lost to starvation, calcium deficiency, or disturbance from other large birds. Cattle egrets maintain a special relationship with cattle, which extends to other large grazing mammals; wider human farming is believed to be a major cause of their suddenly expanded range. The cattle egret removes ticks and flies from cattle and consumes them. This benefits both organisms, but it has been implicated in the spread of tick-borne animal diseases.
Taxonomy
Before the description of the Bubulcus by Charles Lucien Bonaparte in 1855, the western cattle egret had already been described in 1758 by Carl Linnaeus in his Systema Naturae as Ardea ibis, and the eastern cattle egret had been described in 1783 by Pieter Boddaert as Cancroma coromanda. Their generic name Bubulcus is Latin for herdsman, referring, like the English name, to their association with cattle. Ibis is a Latin and Greek word which originally referred to another white wading bird, the sacred ibis, but was applied to the western cattle egret in error. The epithet coromanda refers to the Coromandel Coast of India.
The eastern and western cattle egrets were split by McAllan and Bruce, but were regarded as conspecific by almost all other recent authors until the publication of the influential Birds of South Asia. The eastern cattle egret breeds in South Asia, Eastern Asia, and Australasia, and the western species occupies the rest of the cattle egret's range, including Western Asia, Europe, Africa, and the Americas. According to the IOC birdlist, they are both monotypic species. While some authorities recognise a third Seychelles subspecies, the Seychelles cattle egret (A. i. seychellarum), which was first described by Finn Salomonsen in 1934.
Despite superficial similarities in appearance, the cattle egret is more closely related to the other members of the genus Ardea, which comprises the great or typical herons and the great egret (A. alba), than to the majority of species termed egrets in the genus Egretta. Rare cases of hybridization with little blue herons (Egretta caerulea), little egrets (E. garzetta), and snowy egrets (E. thula) have been recorded.
An older English name for the cattle egret is buff-backed heron.
Description
The cattle egret is a stocky heron with an wingspan; it is long and weighs . It has a relatively short, thick neck, a sturdy bill, and a hunched posture. The nonbreeding adult has mainly white plumage, a yellow bill, and greyish-yellow legs. During the breeding season, adults of the western cattle egret develop orange-buff plumes on the back, breast, and crown, and the bill, legs, and irises become bright red for a brief period prior to pairing. The sexes are similar, but the male is marginally larger and has slightly longer breeding plumes than the female; juvenile birds lack coloured plumes and have a black bill.
The eastern differs from the western in breeding plumage, when the buff colour on its head extends to the cheeks and throat, and the plumes are more golden in colour. This species' bill and tarsi are longer on average than in A. ibis. A. i. seychellarum, which may or may not be a valid subspecies, is smaller and shorter-winged than the other forms. It has white cheeks and throat, like A. ibis, but the nuptial plumes are golden, as with A. coromanda. Individuals with abnormally grey, melanistic plumages have been recorded.
The positioning of the egret's eyes allows for binocular vision during feeding, and physiological studies suggest that they may be capable of crepuscular or nocturnal activity. Adapted to foraging on land, they have lost the ability possessed by their wetland relatives to accurately correct for light refraction by water.
Distribution and habitat
The western cattle egret has undergone one of the most rapid and wide-reaching natural expansions of any bird species. It was originally native to parts of southern Spain and Portugal, tropical and subtropical Africa, and humid tropical and subtropical Asia. At the end of the 19th century, it began expanding its range into southern Africa, first breeding in the Cape Province in 1908. Cattle egrets were first sighted in the Americas on the boundary of Guiana and Suriname in 1877, having apparently flown across the Atlantic Ocean. In the 1930s, the species is thought to have become established in that area. It is now widely distributed across Brazil and was first discovered in the northern region of the country in 1964, feeding along with buffalos.
The species first arrived in North America in 1941 (these early sightings were originally dismissed as escapees), bred in Florida in 1953, and spread rapidly, breeding for the first time in Canada in 1962. It is now commonly seen as far west as California. It was first recorded breeding in Cuba in 1957, in Costa Rica in 1958, and in Mexico in 1963, although it was probably established before then. In Europe, the species had historically declined in Spain and Portugal, but in the latter part of the 20th century, it expanded back through the Iberian Peninsula, and then began to colonise other parts of Europe, southern France in 1958, northern France in 1981, and Italy in 1985. Breeding in the United Kingdom was recorded for the first time in 2008, only a year after an influx seen in the previous year. In 2008, cattle egrets were also reported as having moved into Ireland for the first time. This trend has continued and cattle egrets have become more numerous in southern Britain with influxes in some numbers during the nonbreeding seasons of 2007/08 and 2016/17. They bred in Britain again in 2017, following an influx in the previous winter, and may become established there.
In Australia, the colonisation began in the 1940s, with the eastern cattle egret establishing itself in the north and east of the continent. It began to regularly visit New Zealand in the 1960s. Since 1948, the cattle egret has been permanently resident in Israel. Prior to 1948, it was only a winter visitor.
The massive and rapid expansion of the cattle egret's range is due to its relationship with humans and their domesticated animals. Originally adapted to a commensal relationship with large grazing and browsing animals, it was easily able to switch to domesticated cattle and horses. As the keeping of livestock spread throughout the world, the cattle egret was able to occupy otherwise empty niches. Many populations of cattle egrets are highly migratory and dispersive, and this has helped the genus' range expansion. The cattle egret has been seen as a vagrant in various sub-Antarctic islands, including South Georgia, Marion Island, the South Sandwich Islands, and the South Orkney Islands. A small flock of eight birds was also seen in Fiji in 2008.
In addition to the natural expansion of its range, cattle egrets have been deliberately introduced into a few areas. The western cattle egret was introduced to Hawaii in 1959, and to the Chagos Archipelago in 1955. Successful releases were also made in the Seychelles and Rodrigues, but attempts to introduce them to Mauritius failed. Numerous birds were also released by Whipsnade Zoo in England, but they were never established.
Although the cattle egret sometimes feeds in shallow water, unlike most herons, it is typically found in fields and dry grassy habitats, reflecting its greater dietary reliance on terrestrial insects rather than aquatic prey.
Migration and movements
Some populations of cattle egrets are migratory, others are dispersive, and distinguishing between the two can be difficult. In many areas, populations can be both sedentary and migratory. In the Northern Hemisphere, migration is from cooler climes to warmer areas, but cattle egrets nesting in Australia migrate to cooler Tasmania and New Zealand in the winter and return in the spring. Migration in western Africa is in response to rainfall, and in South America, migrating birds travel south of their breeding range in the nonbreeding season. Populations in southern India appear to show local migrations in response to the monsoons. They move north from Kerala after September. During winter, many birds have been seen flying at night with flocks of Indian pond herons (Ardeola grayii) on the south-eastern coast of India and a winter influx has also been noted in Sri Lanka.
Young birds are known to disperse up to from their breeding area. Flocks may fly vast distances and have been seen over seas and oceans including in the middle of the Atlantic.
Ecology and behavior
Voice
The cattle egret gives a quiet, throaty rick-rack call at the breeding colony, but is otherwise largely silent.
Breeding
The cattle egret nests in colonies, which are often found around bodies of water. The colonies are usually found in woodlands near lakes or rivers, in swamps, or on small inland or coastal islands, and are sometimes shared with other wetland birds, such as herons, egrets, ibises, and cormorants. The breeding season varies within South Asia. Nesting in northern India begins with the onset of monsoons in May. The breeding season in Australia is November to early January, with one brood laid per season. The North American breeding season lasts from April to October. In the Seychelles, the breeding season of B. i. seychellarum is April to October.
The male displays in a tree in the colony, using a range of ritualised behaviours, such as shaking a twig and sky-pointing (raising his bill vertically upwards), and the pair forms over 3–4 days. A new mate is chosen in each season and when renesting following nest failure. The nest is a small, untidy platform of sticks in a tree or shrub constructed by both parents. Sticks are collected by the male and arranged by the female, and stick-stealing is rife. The clutch size can be one to five eggs, although three or four is most common. The pale bluish-white eggs are oval-shaped and measure . Incubation lasts around 23 days, with both sexes sharing incubation duties. The chicks are partly covered with down at hatching, but are not capable of fending for themselves; they become capable of regulating their temperature at 9–12 days and are fully feathered in 13–21 days. They begin to leave the nest and climb around at 2 weeks, fledge at 30 days and become independent at around the 45th day.
The cattle egret engages in low levels of brood parasitism, and a few instances have been reported of cattle egret eggs being laid in the nests of snowy egrets and little blue herons, although these eggs seldom hatch. Also, evidence of low levels of intraspecific brood parasitism has been found, with females laying eggs in the nests of other cattle egrets. As much as 30% extra-pair copulations has been noted.
The dominant factor in nesting mortality is starvation. Sibling rivalry can be intense, and in South Africa, third and fourth chicks inevitably starve. In the dryer habitats with fewer amphibians, the diet may lack sufficient vertebrate content and may cause bone abnormalities in growing chicks due to calcium deficiency. In Barbados, nests were sometimes raided by vervet monkeys, and a study in Florida reported the fish crow and black rat as other possible nest raiders. The same study attributed some nestling mortality to brown pelicans nesting in the vicinity, which accidentally, but frequently, dislodged nests or caused nestlings to fall. In Australia, Torresian crows, wedge-tailed eagles, and white-bellied sea eagles take eggs or young, and tick infestation and viral infections may also be causes of mortality.
Feeding
The cattle egret feeds on a wide range of prey, particularly insects, especially grasshoppers, crickets, flies (adults and maggots), beetles, and moths, as well as spiders, frogs, fish, crayfish, small snakes, lizards and earthworms. In a rare instance, they have been observed foraging along the branches of a banyan tree for ripe figs. The cattle egret is usually found with cattle and other large grazing and browsing animals, and catches small creatures disturbed by the mammals. Studies have shown that cattle egret foraging success is much higher when foraging near a large animal than when feeding singly. When foraging with cattle, it has been shown to be 3.6 times more successful in capturing prey than when foraging alone. Its performance is similar when it follows farm machinery, but it is forced to move more. In urban situations, cattle egrets have also been observed foraging in peculiar situations such as railway lines.
A cattle egret will weakly defend the area around a grazing animal against others of the same species, but if the area is swamped by egrets, it will give up and continue foraging elsewhere. Where numerous large animals are present, cattle egrets selectively forage around species that move at around 5–15 steps per minute, avoiding faster and slower moving herds; in Africa, cattle egrets selectively forage behind plains zebras, waterbuck, blue wildebeest and Cape buffalo. Dominant birds feed nearest to the host, and thus obtain more food.
The cattle egret sometimes shows versatility in its diet. On islands with seabird colonies, it will prey on the eggs and chicks of terns and other seabirds. During migration, it has also been reported to eat exhausted migrating landbirds. Birds of the Seychelles race also indulge in some kleptoparasitism, chasing the chicks of sooty terns and forcing them to disgorge food.
Threats
Pairs of crested caracaras have been observed chasing cattle egrets in flight, forcing them to the ground, and killing them.
Status
The IUCN Red List treats them as a single species. They have a large range, with an estimated global extent of occurrence of . Their global population is estimated to be 3.8–6.7 million individuals. For these reasons, the genus is evaluated as least concern. The expansion and establishment of the genus over large ranges has led it to be classed as an invasive species, although little, if any, impact has been noted yet.
Relationship with humans
As a conspicuous genus, the cattle egret has attracted many common names. These mostly relate to its habit of following cattle and other large animals, and it is known variously as cow crane, cow bird or cow heron, or even elephant bird or rhinoceros egret. Its Arabic name, abu qerdan, means "father of ticks", a name derived from the huge number of parasites such as avian ticks found in its breeding colonies. The Maasai people consider the presence of large numbers of cattle egrets as an indicator of impending drought and use it to decide on moving their cattle herds.
Cattle egrets are an occurring traditional motif in fishing boats among fishermen of the Malay Peninsula east coast who believed them as a symbol of good luck and fortune.
The cattle egret is a popular bird with cattle ranchers for its perceived role as a biocontrol of cattle parasites such as ticks and flies. A study in Australia found that cattle egrets reduced the number of flies that bothered cattle by pecking them directly off the skin. It was the benefit to stock that prompted ranchers and the Hawaiian Board of Agriculture and Forestry to release the western cattle egret in Hawaii.
Not all interactions between humans and cattle egrets are beneficial. The cattle egret can be a safety hazard to aircraft due to its habit of feeding in large groups in the grassy verges of airports, and it has been implicated in the spread of animal infections such as heartwater, infectious bursal disease, and possibly Newcastle disease.
| Biology and health sciences | Pelecanimorphae | Animals |
319558 | https://en.wikipedia.org/wiki/Ocean%20liner | Ocean liner | An ocean liner is a type of passenger ship primarily used for transportation across seas or oceans. Ocean liners may also carry cargo or mail, and may sometimes be used for other purposes (such as for pleasure cruises or as hospital ships). The Queen Mary 2 is the only ocean liner still in service to this day, serving with Cunard Line.
The category does not include ferries or other vessels engaged in short-sea trading, nor dedicated cruise ships where the voyage itself, and not transportation, is the primary purpose of the trip. Nor does it include tramp steamers, even those equipped to handle limited numbers of passengers. Some shipping companies refer to themselves as "lines" and their passenger ships, which often operate over set routes according to established schedules, as "liners".
Though ocean liners share certain similarities with cruise ships, they must be able to travel between continents from point A to point B on a fixed schedule, so must be faster and built to withstand the rough seas and adverse conditions encountered on long voyages across the open ocean. To protect against large waves they usually have a higher hull and promenade deck with higher positioning of lifeboats (the height above water called the freeboard), as well as a longer bow than a cruise ship. Additionally, for additional strength they are often designed with thicker hull plating than is found on cruise ships, as well as a deeper draft for greater stability, and have large capacities for fuel, food, and other consumables on long voyages. On an ocean liner, the captain's tower (bridge) is usually positioned on the upper deck for increased visibility.
The first ocean liners were built in the mid-19th century. Technological innovations such as the steam engine, Diesel engine and steel hull allowed larger and faster liners to be built, giving rise to a competition between world powers of the time, especially between the United Kingdom, the German Empire, and to a lesser extent France. Once the dominant form of travel between continents, ocean liners were rendered largely obsolete by the emergence of long-distance aircraft after World War II. Advances in automobile and railway technology also played a role. After was retired in 2008, the only ship still in service as an ocean liner is , introduced in 2004, as well as the largest ever built.
Overview
Ocean liners were the primary mode of intercontinental travel for over a century, from the mid-19th century until they began to be supplanted by airliners in the 1950s. In addition to passengers, liners carried mail and cargo. Ships contracted to carry British Royal Mail used the designation RMS. Liners were also the preferred way to move gold and other high-value cargoes.
The busiest route for liners was on the North Atlantic with ships travelling between Europe and North America. It was on this route that the fastest, largest and most advanced liners travelled, though most ocean liners historically were mid-sized vessels which served as the common carriers of passengers and freight between nations and among other countries and their colonies and dependencies before the dawn of the jet age. Such routes included Europe to African and Asian colonies, Europe to South America, and migrant traffic from Europe to North America in the 19th and first two decades of the 20th centuries, and to Canada and Australia after the Second World War.
Shipping lines are companies engaged in shipping passengers and cargo, often on established routes and schedules. Regular scheduled voyages on a set route are called "line voyages" and vessels (passenger or cargo) trading on these routes to a timetable are called liners. The alternative to liner trade is "tramping" whereby vessels are notified on an ad hoc basis as to the availability of a cargo to be transported. (In older usage, "liner" also referred to ships of the line, that is, line-of-battle ships, but that usage is now rare.) The term "ocean liner" has come to be used interchangeably with "passenger liner", although it can refer to a cargo liner or cargo-passenger liner.
The advent of the Jet Age and the decline in transoceanic ship service brought about a gradual transition from passenger ships to modern cruise ships as a means of transportation. In order for ocean liners to remain profitable, cruise lines modified some of them to operate on cruise routes, such as the . Certain characteristics of older ocean liners made them unsuitable for cruising, such as high fuel consumption, deep draught preventing them from entering shallow ports, and cabins (often windowless) designed to maximize passenger numbers rather than comfort. The Italian Line's and , the last ocean liners to be built primarily for crossing the North Atlantic, could not be converted economically and had short careers.
History
19th century
At the beginning of the 19th century, the Industrial Revolution and the inter-continental trade made the development of secure links between continents imperative. Being at the top among the colonial powers, the United Kingdom needed stable maritime routes to connect different parts of its empire: the Far East, India, Australia, etc. The birth of the concept of international water and the lack of any claim to it simplified navigation during this period. In 1818, the Black Ball Line, with a fleet of sailing ships, offered the first regular passenger service with emphasis on passenger comfort, from England to the United States.
In 1807, Robert Fulton succeeded in applying steam engines to ships. He built the first ship that was powered by this technology, the Clermont, which succeeded in travelling between New York City and Albany, New York in thirty hours before entering into regular service between the two cities. Soon after, other ships were built using this innovation. In 1816, the became the first steamship to cross the English Channel. Another important advance came in 1819, when became the first steamship to cross the Atlantic Ocean. She left the U.S. city of the same name and arrived in Liverpool, England in 27 days. Most of the distance was covered by sailing; the steam power was not used for more than 72 hours during the travel. Public enthusiasm for the new technology was not high, as none of the thirty-two people who had booked a seat boarded the ship for that historic voyage. Although Savannah had proven that a steamship was capable of crossing the ocean, the public was not yet prepared to trust such means of travel on an open sea, and, in 1820, the steam engine was removed from the ship.
Work on this technology continued and a new step was taken in 1833. Royal William managed to cross the Atlantic by using steam power on most of the voyage; sail was used only when the boilers were cleaned. However, there were still many skeptics, and in 1836, scientific writer Dionysius Lardner declared that: As the project of making the voyage directly from New York to Liverpool, it was perfectly chimerical, and they might as well talk of making the voyage from New York to the moon.
The last step toward long-distance travel using steam power was taken in 1837 when left Liverpool on 4 April and arrived in New York eighteen days later on 22 April after a turbulent crossing. Too little coal was prepared for the crossing, and the crew had to burn cabin furniture in order to complete the voyage. The journey took place at a speed of 8.03 knots. The voyage was made possible by the use of a condenser, which fed the boilers with fresh water and avoiding having to periodically shut down the boilers in order to remove the salt. This new record was short-lived. The next day, , designed by railway engineer Isambard Kingdom Brunel, arrived in New York. She left Liverpool on 8 April and overtook Siriuss record with an average speed of 8.66 knots. A race for speed was born, and, with it, the tradition of the Blue Riband.
With Great Western, Isambard Kingdom Brunel laid the foundations for new shipbuilding techniques. He realised that the carrying capacity of a ship increases as the cube of its dimensions, whilst the water resistance only increases as the square of its dimensions. This means that large ships are more fuel-efficient, something very important for long voyages across the Atlantic. Constructing large ships was therefore more profitable. Moreover, migration to the Americas increased enormously. These movements of population were a financial windfall for the shipping companies, of which some of the largest were founded during this period. Examples are the P&O of the United Kingdom in 1822 and the Compagnie Générale Transatlantique of France in 1855.
The steam engine also allowed ships to provide regular service without the use of sail. This aspect particularly appealed to the postal companies, which leased the services of ships to serve clients separated by the ocean. In 1839, Samuel Cunard founded the Cunard Line and became the first to dedicate the activity of his shipping company to the transport of mails, thus ensuring regular services on a given schedule. The company's ships operated the routes between the United Kingdom and the United States. Over time, the paddle wheel, impractical on the high seas, was abandoned in favour of the propeller. In 1840, Cunard Line's began its first regular passenger and cargo service by a steamship, sailing from Liverpool to Boston, Massachusetts.
As the size of ships increased, the wooden hull became unreliable. The start of the use of iron hulls in 1845, and then of steel hulls, solved this problem. The first ship to be both iron-hulled and equipped with a screw propeller was , a creation of Brunel. Her career was disastrous and short. She was run aground and stranded at Dundrum Bay in 1846. In 1884, she was retired to the Falkland Islands where she was used as a warehouse, quarantine ship, and coal hulk until she was scuttled in 1937. The American company Collins Line took a different approach. It equipped its ships with cold rooms, heating systems, and various other innovations but the operation was expensive. The sinking of two of its ships was a major blow to the company which was dissolved in 1858.
In 1858, Brunel built his third and last giant, . The ship was, for 43 years, the largest passenger ship ever built. She had the capacity to carry 4,000 passengers. Her career was marked by a series of failures and incidents, one of which was an explosion on board during her maiden voyage.
Many ships owned by German companies such as Hamburg America Line and Norddeutscher Lloyd were sailing from major German ports, such as Hamburg and Bremen, to the United States during this time. The year 1858 was marked by a major accident: the sinking of . The ship, built in Greenock and sailing between Hamburg and New York twice a month, suffered an accidental fire off the coast of Newfoundland and sank with the loss of all but 89 of the 542 passengers.
In the British market, Cunard Line and White Star Line competed strongly against each other in the late 1860s. The struggle was symbolised by the attainment of the Blue Riband, which the two companies achieved several times around the end of the century. The luxury and technology of ships were also evolving. Auxiliary sails became obsolete and disappeared completely at the end of the century. Possible military use of passenger ships was envisaged and, in 1889, became the first auxiliary cruiser in history. In the time of war, ships could easily be equipped with cannons and used in cases of conflict. Teutonic succeeded in impressing Emperor Wilhelm II of Germany, who wanted to see his country endowed with a modern fleet.
In 1870, the White Star Line's set a new standard for ocean travel by having its first-class cabins amidships, with the added amenity of large portholes, electricity and running water. The size of ocean liners increased from 1880 onward to meet the needs of migration to the United States and Australia.
and her sister ship were the last two Cunard liners of the period to be fitted with auxiliary sails. Both ships were built by John Elder & Co. of Glasgow, Scotland, in 1884. They were record breakers by the standards of the time, and were the largest liners then in service, serving the Liverpool to New York route.
was a 6,814-ton steamship owned by the Orient Steamship Co., and was fitted with refrigeration equipment. She served the Suez Canal route from England to Australia during the 1890s, up until the years leading to World War I when she was converted to an armed merchant cruiser.
In 1897, Norddeutscher Lloyd launched . She was followed three years later by three sister ships. The ship was both luxurious and fast, managing to win the Blue Riband from the British. She was also the first of the fourteen ocean liners with four funnels that have emerged in maritime history. The ship needed only two funnels, but more funnels gave passengers a feeling of safety. In 1900, the Hamburg America Line competed with its own four-funnel liner, . She quickly obtained the Blue Riband for her company. This race for speed, however, was a detriment to passengers' comfort and generated strong vibration, which made her owner lose any interest in her after she lost the Blue Riband to another ship of Norddeutscher Lloyd. She was only used for ten years for transatlantic crossing before being converted into a cruise ship. Until 1907, the Blue Riband remained in the hands of the Germans.
Early 20th century
In 1902, J. P. Morgan embraced the idea of a maritime empire comprising a large number of companies. He founded the International Mercantile Marine Co., a trust which originally comprised only American shipping companies. The trust then absorbed Leyland Line and White Star Line. The British government then decided to intervene in order to regain its ascendancy.
Although German liners dominated in terms of speed, British liners dominated in terms of size. and the Big Four of the White Star Line were the first liners to surpass Great Eastern as the largest passenger ships. Ultimately their owner was American (as mentioned above, White Star Line had been absorbed into J. P. Morgan's trust). Faced with this major competition, the British government contributed financially to Cunard Line's construction of two liners of unmatched size and speed, under the condition that they be available for conversion into armed cruisers when needed by the navy. The result of this partnership was the completion in 1907 of two sister ships: and , both of which won the Blue Riband during their respective maiden voyages. The latter retained it for twenty years. Their great speed was achieved by the use of turbines instead of conventional expansion machines.
In response to the competition from Cunard Line, White Star Line ordered the liners at the end of 1907. The first of these three liners, , completed in 1911, had a fine career, although punctuated by incidents. This was not the case for her sister, the , which sank on her maiden voyage on 15 April 1912, resulting in several major changes to maritime safety practices. As for the third sister, , she never served her intended purpose as a passenger ship, as she was drafted in the First World War as a hospital ship, and sank to a naval mine in 1916.
At the same time, France tried to mark its presence with the completion in 1912 of owned by the Compagnie Générale Transatlantique. Germany soon responded to the competition from the British. From 1912 to 1914, Hamburg America Line completed a trio of liners significantly larger than the White Star Line's Olympic-class ships. The first to be completed, in 1913, was . She was followed by SS Vaterland in 1914. The construction of the third liner, , was paused by the outbreak of World War I.
World War I was a difficult time for the liners. Some of them, like the Mauretania, , and Britannic were transformed into hospital ships during the conflict. Others became troop transports, while some, such as the Kaiser Wilhelm der Grosse, participated in the war as warships. Troop transportation was very popular due to the liners' large size. Liners converted into troop ships were painted in dazzle camouflage to reduce the risk of being torpedoed by enemy submarines.
The war caused the loss of many liners. Britannic, while serving as a hospital ship, sank in the Aegean Sea in 1916 after she struck a mine. Numerous incidents of torpedoing took place and large numbers of ships sank. Kaiser Wilhelm der Grosse was defeated and scuttled after a fierce battle with off the coast of west Africa, while her sister ship served as a commerce raider. The torpedoing and sinking of Lusitania on 7 May 1915 caused the loss of 128 American lives at a time when the United States was still neutral. Although other factors came into play, the loss of American lives in the sinking strongly pushed the United States to favour the Allied Powers and facilitated the country's entry into the war.
The losses of the liners owned by the Allied Powers were compensated by the Treaty of Versailles in 1919. This led to the awarding of many German liners to the victorious Allies. The Hamburg America Line's trio (, Vaterland, and Bismarck) were divided between the Cunard Line, White Star Line, and the United States Lines, while the three surviving ships of the Kaiser class were requisitioned by the US Navy in the context of the conflict and then retained. The Tirpitz, whose construction was delayed by the outbreak of war, eventually became the . Of the German superliners, only Deutschland, because of her poor state, avoided this fate.
After World War I
After a period of reconstruction, the shipping companies recovered quickly from the damage caused by World War I. The ships whose construction was started before the war, such as of the French Line, were completed and put into service. Prominent British liners, such as the Olympic and the Mauretania, were also put back into service and had a successful career in the early 1920s. More modern liners were also built, such as (completed in 1927). The United States Lines, having received the Vaterland, renamed her Leviathan and made her the flagship of the company's fleet. Because all U.S. registered ships counted as an extension of U.S. territory, the National Prohibition Act made American liners alcohol-free, causing alcohol-seeking passengers to choose ships of other countries for travel and substantially reduce profits for the United States Lines.
In 1929, Germany returned to the scene with the two ships of Norddeutscher Lloyd, and . Bremen won the Blue Riband from Britain's Mauretania after the latter had held it for twenty years. Soon, Italy also entered the scene. The Italian Line completed and in 1932, breaking the records of both luxury and speed (Rex won the westbound Blue Riband in 1933). France reentered the scene with of the French Compagnie Générale Transatlantique (CGT). The ship was the largest ship afloat at the time of her completion in 1935. She was also the fastest, winning the Blue Riband in 1935.
A crisis arose when the United States drastically reduced its immigrant quotas, causing shipping companies to lose a large part of their income and to have to adapt to this circumstance. The Great Depression also played an important role, causing a drastic decrease in the number of people crossing the Atlantic and at the same time reducing the number of profitable transatlantic voyages. In response, shipping companies redirected many of their liners to a more profitable cruise service. In 1934, in the United Kingdom, Cunard Line and White Star Line were in very bad shape financially. Chancellor of the Exchequer Neville Chamberlain proposed to merge the two companies in order to solve their financial problems. The merger took place in 1934 and launched the construction of the while progressively sending their older ships to the scrapyard. The Queen Mary was the fastest ship of her time and the largest for a short amount of time, she captured the Blue Riband twice, both off Normandie. The construction of a second ship, the , was interrupted by the outbreak of World War II.
World War II was a conflict rich in events involving liners. From the start of the conflict, German liners were requisitioned and many were turned into barracks ships. It was in the course of this activity that the Bremen caught fire while under conversion for Operation Sea Lion and was scrapped in 1941. During the conflict, Queen Elizabeth and Queen Mary provided distinguished service as troopships.
Many liners were sunk with great loss of life; the three worst disasters were the loss of the Cunarder in 1940 off Saint-Nazaire to German bombing while attempting to evacuate troops of the British Expeditionary Force from France, with the loss of more than 3,000 lives; the sinking of , after the ship was torpedoed by a Soviet submarine, with more than 9,000 lives lost, making it the deadliest maritime disaster in history; and the sinking of with more than 7,000 lives lost, both in the Baltic Sea, in 1945.
SS Rex was bombarded and sunk in 1944, and Normandie caught fire, capsized, and sank in New York in 1942 while being converted into a troopship. Many of the superliners of the 1920s and 1930s were victims of U-boats, mines or enemy aircraft. was attacked by German planes, then torpedoed by a U-boat when tugs tried to tow her to safety. Out of all the innovative and glamorous inter-war superliners, only the Cunard Queens and Europa would survive the war.
Decline of long-distance line voyages
After the war, some ships were again transferred from the defeated nations to the winning nations as war reparations. This was the case of the Europa, which was ceded to France and renamed Liberté. The United States government was very impressed with the service of the Cunard's Queen Mary and Queen Elizabeth as troopships during the war. To ensure a reliable and fast troop transport in case of a war against the Soviet Union, the U.S. government sponsored the construction of and entered it into service for the United States Lines in 1952. She won the Blue Riband on her maiden voyage in that year and held it until Richard Branson won it in 1986 with Virgin Atlantic Challenger II. One year later, in 1953, Italy completed the , which later sank in 1956 after a collision with .
Before World War II, aircraft had not posed a significant economic threat to ocean liners. Most pre-war aircraft were noisy, vulnerable to bad weather, and/or incapable of the range needed for transoceanic flights; all were expensive and had a small passenger capacity. The war accelerated development of large, long-ranged aircraft. Four-engined bombers, such as the Avro Lancaster and Boeing B-29 Superfortress, with their range and massive carrying capacity, were natural prototypes for post-war next-generation airliners. Jet engine technology also accelerated due to wartime development of jet aircraft. In 1953, the De Havilland Comet became the first commercial jet airliner; the Sud Aviation Caravelle, Boeing 707 and Douglas DC-8 followed, and much long-distance travel was done by air. The Italian Line's and , launched in 1962 and 1963, were two of the last ocean liners to be built primarily for liner service across the North Atlantic. Cunard's transatlantic liner, , although designed as an ocean liner, was also used as a cruise ship. By the early 1960s, 95% of passenger traffic across the Atlantic was by aircraft. Thus the reign of the ocean liners came to an end. By the early 1970s, many passenger ships continued their service in cruising.
In 1982, during the Falklands War, three active or former liners were requisitioned for war service by the British Government. The liners Queen Elizabeth 2 and , were requisitioned from Cunard and P&O to serve as troopships, carrying British Army personnel to Ascension Island and the Falkland Islands to recover the Falklands from the invading Argentine forces. The P&O educational cruise ship and former British India Steam Navigation Company liner was requisitioned as a hospital ship, and served after the war as a troopship until the RAF Mount Pleasant station was built at Stanley, which could handle trooping flights.
21st century
By the first decade of the 21st century, only a few former ocean liners were still in existence. Some, like , were sailing as cruise ships while others, like , were preserved as museums, or laid up at pier side like SS United States. After the retirement of Queen Elizabeth 2 in 2008, the only ocean liner in service was Queen Mary 2, built in 2003–04 and used for both point-to-point line voyages and for cruises.
A proposed and planned ocean liner, the Titanic II, is a modern replica of the original RMS Titanic, which sank in 1912. The ship is owned by Blue Star Line and is bought by Australian businessman Clive Palmer. The ship is set to be launched by 2027.
Survivors
Four ocean liners built before the World War II survive today as they have been partially or fully preserved as museums and hotels. The Japanese ocean liner (1929), has been preserved in Naka-ku, Yokohama, Japan, as a museum ship, since 1961. (1934) was preserved in 1967 after her retirement, and became a museum/hotel in Long Beach, California. In the 1970s, (1843) was also preserved, and now resides in Bristol, England as another museum. The latest ship to undergo preservation is (1914). While originally being a cargo ship, it served as the Italian ocean liner Franca C. for Costa Lines from 1952 to 1959, and in 2010 it became a dry berthed luxury hotel on Bintan Island, Indonesia.
Post-war ocean liners still existent include (1948), (1952), MV Brazil Maru (1954), (1958), (1961), (1962), Queen Elizabeth 2 (1967), and Queen Mary 2 (2003). Out of these eight ocean liners, only one is still active and three of them have since been preserved. The Rotterdam has been moored in Rotterdam as a museum and hotel since 2008, while the Queen Elizabeth 2 has been a floating luxury hotel and museum at Mina Rashid, Dubai since 2018. The Ancerville was refurbished as a hotel for use at the Sea World development in Shenzhen, China in 1984. The first of these, Astoria (originally the ocean liner MS Stockholm, which collided with Andrea Doria in 1956) has been rebuilt and refitted as a cruise ship over the years and was in active service for Cruise & Maritime Voyages until operations ceased in 2020 due to the COVID-19 pandemic. In August, 2021 she was purchased by Brock Pierce to be transformed into a hotel along with . These plans were ultimately abandoned and the ship was again made available for sale, never having left port in Rotterdam. Astoria was reported to have been sold for scrap in January 2023, but this has been denied by the ship's owner. United States has been docked in Philadelphia since 1996, but following a legal dispute between the organization that owns United States and the pier owners, she was purchased by Okaloosa County, Florida to be turned into the world's largest artificial reef. There are plans for a land-based museum and several pieces of United States are planned to be preserved. Brazil Maru was beached in Zhanjiang, China as a tourist attraction called Hai Shang Cheng Shi in 1998, though has been closed as of 2022. Funchal was purchased by Brock Pierce in 2021, with the intent of turning her into a hotel. Her future is uncertain as it was reported in July 2021 that no progress has been made since then.
Characteristics
Size and speed
Since their beginning in the 19th century, ocean liners needed to meet growing demands. The first liners were small and overcrowded, leading to unsanitary conditions on board. Eliminating these phenomena required larger ships, to reduce the crowding of passengers, and faster ships, to reduce the duration of transatlantic crossings. The iron and steel hulls and steam power allowed for these to be achieved. Thus, SS Great Western (1,340 GRT) and SS Great Eastern (18,915 GRT) were constructed in 1838 and 1858 respectively. The record set by SS Great Eastern was not beaten until 43 years later in 1901 when (20,904 GT) was completed. The tonnage then grew profoundly: the first liners to have a tonnage that exceeded 20,000 were the Big Four of the White Star Line. The liners, first completed in 1911, were the first to have a tonnage that exceeded 45,000 and the liners first completed in 1913 became the 1st liners with tonnage exceeding 50,000. , completed in 1935, had a tonnage of 79,280. In 1940, raised the record of size to a tonnage of 83,673. She was the largest passenger ship ever constructed until 1997. In 2003, became the largest, at 149,215 GT.
In the early 1840s, the average speed of liners was less than 10 knots (a crossing of the Atlantic thus took about 12 days or more). In the 1870s, the average speed of liners increased to around 15 knots the duration of a transatlantic crossing shortened to around 7 days, owing to the technological progress made in the propulsion of ships. The rudimentary steam boilers gave rise to more elaborate machineries and the paddlewheel gradually disappeared, replaced first by one screw then by two screws. At the beginning of the 20th century, Cunard Line's and reached a speed of 27 knots. Their records seemed unbeatable, and most shipping companies abandoned the race for speed in favor of size, luxury, and safety. The advent of ships with diesel engines, and of those whose engines were oil-burning, such as the Bremen, in the early 1930s, renewed the race for the Blue Riband. The won it in 1935 before being snatched by in 1938. It was not until 1952 that set a record that remains today: 34.5 knots (3 days and 12 hours of crossing the Atlantic). In addition, since 1935, the Blue Riband is accompanied by the Hales Trophy, which is awarded to the winner.
Passenger cabins and amenities
The first ocean liners were designed to carry mostly migrants. On-board sanitary conditions were often deplorable and epidemics were frequent. In 1848, maritime laws imposing hygiene rules were adopted and they improved on-board living conditions. Gradually, two distinct classes were developed: the cabin class and the steerage class. The passengers travelling on the former were wealthy passengers and they enjoyed certain comfort in that class. The passengers travelling on the latter were members of the middle class or the working class. In that class, they were packed in large dormitories. Until the beginning of the 20th century, they did not always have bedsheets and meals. An intermediate class for tourists and members of the middle class gradually appeared. The cabins were then divided into three classes.
The facilities offered to passengers developed over time. In the 1870s, the installation of bathtubs and oil lamps caused a sensation on board . In the following years, the number of amenities became numerous, for example: smoking rooms, lounges, and promenade deck. In 1907, even offered Turkish baths and a swimming pool. In the 1920s, was the first liner to offer a movie theatre.
Builders
British and German
The British and the German shipyards were the most famed in shipbuilding during the era of ocean liners. In Ireland, Harland & Wolff shipyard of Belfast were particularly innovative and succeeded in winning the trust of many shipping companies, such as White Star Line. These gigantic shipyards employed a large portion of the population of cities and built hulls, machines, furnitures and lifeboats. Among the other well-known British shipyards were Swan, Hunter & Wigham Richardson, the builder of , and John Brown & Company, builders of , , , , and Queen Elizabeth 2.
Germany had many shipyards on the coast of the North Sea and the Baltic Sea, including Blohm & Voss and AG Vulcan Stettin. Many of these shipyards were destroyed during World War II; some managed to recover and continue building ships.
Other nations
In France, major shipyards included Chantiers de Penhoët in Saint-Nazaire, known for building . This shipyard merged with Ateliers et Chantiers de la Loire shipyard to form the Chantiers de l'Atlantique shipyard, which has built ships including . France also had major shipyards on the shores of the Mediterranean Sea.
Italy and the Netherlands also had shipyards capable of building large ships (for example, Fincantieri).
Shipping companies
British
There were many British shipping companies; two were particularly distinguished: Cunard Line and White Star Line. Both were founded during the 1840s and engaged in strong competition against one another, possessing the largest and fastest liners in the world in the early 20th century. It was not until 1934 that financial difficulty caused the two to merge, forming Cunard White Star Ltd. The P&O also occupied a large part of the business.
The Royal Mail Steam Packet Company operated as a state-owned enterprise with its close relationship with the government. Over the course of its history, it took over many shipping companies, becoming one of the largest companies in the world before legal problems led to its liquidation in 1931. The Union Castle Line operated in Africa and the Indian Ocean with a fleet of considerable size.
German, French and Dutch
Two rival companies, Hamburg America Line (often referred to as "HAPAG") and Norddeutscher Lloyd, competed in Germany. The First and Second World Wars dealt much damage to the two companies, both forced to cede their ships to the winning side in both wars. The two merged to form Hapag-Lloyd in 1970.
The ocean liner industry in France also consisted of two rival companies: the Compagnie Générale Transatlantique (commonly known as "Transat" or "French Line") and Messageries Maritimes. The CGT operated on the North Atlantic route with well-known liners such as and , while the MM operated in French colonies in Asia and Africa. Decolonization in the second half of the 20th century led to a sharp decline in profit for the MM, and it merged with the CGT in 1975 to form the Compagnie Générale Maritime.
The Netherlands had three main companies. The Holland America Line operated mostly on the north Atlantic route and with well-known ships like the and . Unlike the French and German industry, the Holland America Line had no domestic rival in this trade and only had to compete with foreign lines. The other two Dutch lines were the Stoomvaart Maatschappij Nederland (SMN), otherwise known as the Netherland Line and the Koninklijke Rotterdamsche Lloyd (KRL); both offered regular service between the Netherlands and the Dutch East Indies, the Dutch colony in South East Asia now known as Indonesia, and had a long-lasting friendly rivalry.
Other nations
The United States Lines competed with European companies for the North Atlantic trade. In Italy, the Italian Line was founded in 1932 as a result of a merger of three companies. It was known for operating liners such as and . The Japanese established Nippon Yusen, also known as NYK Lines, which ran trans-Pacific liners such as the Hikawa Maru and the Asama Maru.
Routes
North Atlantic
The most important of all routes taken by ocean liners was the North Atlantic route. It accounted for a large part of the clientele, who traveled between ports of Liverpool, Southampton, Hamburg, Le Havre, Cherbourg, Cobh, and New York City. The profitability of this route came from migration to the United States. The need for speed influenced the construction of liners for this route, and the Blue Riband was awarded to the liner with the highest speed. The route was not without danger, as storm and icebergs are common in the North Atlantic. Many shipwrecks occurred on this route, among them that of , the details of which have been recounted in numerous books, films and documentaries. This route was the preferred route for major shipping companies and was the scene of fierce competition between them.
South Atlantic
The South Atlantic was the route frequented by liners bound for South America, Africa, and sometimes Oceania. The White Star Line had some of its ships, such as the , on the Liverpool-Cape Town-Sydney route. There was not the same level of competition in the South Atlantic as there was in the North Atlantic. There were fewer shipwrecks. The Hamburg Süd operated on this route; among its ships was the famed .
Mediterranean
The Mediterranean Sea was frequented by many ocean liners. Many companies benefited from migration from Italy and the Balkans to the United States. Cunard's served on the Gibraltar-Genoa-Trieste route. Similarly, Italian liners crossed the Mediterranean Sea before entering the North Atlantic Ocean. The opening of the Suez Canal made the Mediterranean a possible route to Asia.
Indian Ocean and the East Asia
Colonization made Asia particularly attractive to shipping companies. Many government officials must travel there from time to time. As early as the 1840s, the P&O organized trips to Calcutta via the Suez Isthmus, as the canal had not yet been built. The time it took to travel on this route to India, Southeast Asia, and Japan was long, with many stopovers. The Messageries Maritimes operated on this route, notably in the 1930s, with its motor ships. Similarly, the La Marseillaise, put into service in 1949, was one of the flagships of its fleet. Decolonization caused the loss in the profitability of these ships.
Pacific
Ocean liners on the Pacific route brought large numbers of migrants from East Asia to the Americas, especially the United States, which continued despite successive laws restricting Asian immigration to the United States; the journey typically took three weeks, with many impoverished migrants travelling in steerage class conditions. Some of the finest ships on the route, such as of Canadian Pacific Steamships which operated out of Vancouver, and Hikawa Maru of Nippon Yusen, became known as 'Queen of the Pacific'.
Other
National symbol
The construction of some ocean liners was a result of nationalism. The revival of power of the German navy stemmed from the clear affirmation of Kaiser Wilhelm II of Germany to see his country become a sea power. Thus, the of 1900 had the honor to bear the name of its mother country, an honor which she lost after ten years of a disappointing career. and of 1907 were built with the help of the British government with the desire that the United Kingdom would regain its prestige as a sea power. of 1952 was the result of a desire by the United States government to possess a large and fast ship that is convertible into a troop transport. and of 1932 were constructed at the demands of Benito Mussolini. Finally, the construction of 1961 was a result of Charles de Gaulle's desire to build on French national pride and was financed by the French government.
Some liners did gain great popularity. Mauretania and had many admirers during their careers, and their retirement and scrapping caused some sadness. The same was true of Île de France, whose scrapping aroused strong emotion from her admirers. Similarly, was very popular with the British people.
Maritime disasters and incidents
Some ocean liners are known today because of their sinking with great loss of lives. In 1873 struck an underwater rock and sank off the coast of Nova Scotia, Canada, killing at least 535 people. In 1912 the sinking of the RMS Titanic, which took approximately 1,500 lives, highlighted the overconfidence of the shipping companies in their ships, such as the failure to put enough lifeboats on board. Safety measures at sea were reexamined following the incident. Two years later, in 1914, sank in the Saint Lawrence River after colliding with the . 1,012 people died.
Among the other sinkings are the sinking by torpedo of the RMS Lusitania in 1915, which resulted in the loss of 1,198 lives and provoked an international outcry, the sinking by naval mine of in 1916, and that of , which caught fire and sank in the Gulf of Aden in 1932, killing 54 people. In 1956 the sinking of , with the loss of 46 lives, after a collision with made the headline.
In 1985, was hijacked off the coast of Egypt by members of the Palestinian Liberation Front, resulting in the death of one of the hostages being held by the hijackers. In 1994, she caught fire and sank off the coast of Somalia.
In popular culture
Literature
Ocean liners have a strong impact on popular culture, whether during their era or afterwards. In 1867, Jules Verne recounted his experience aboard in his novel A Floating City. In 1898, writer Morgan Robertson wrote the short novel Futility, or the Wreck of the Titan, which features a British ocean liner Titan that hits an iceberg and sinks in the North Atlantic with great loss of lives. The similarities between the plot of the novel and the sinking of the 14 years later led to the assertion of conspiracy theories regarding Titanic.
Films
Ocean liners were often a setting of a love story in films, such as the 1939's Love Affair Liners were also used as a setting of disaster films. The 1960 film The Last Voyage was filmed on board the Île de France, which was used as a floating prop and was scuttled for the occasion. The 1972 film The Poseidon Adventure has become a classic of the genre and has spawned many remakes.
The sinking of Titanic also attracted attention of filmmakers. Nearly fifteen films were made to depict it, with James Cameron's 1997 film being the most well-known and commercially successful.
| Technology | Naval transport | null |
319610 | https://en.wikipedia.org/wiki/Biological%20life%20cycle | Biological life cycle | In biology, a biological life cycle (or just life cycle when the biological context is clear) is a series of stages of the life of an organism, that begins as a zygote, often in an egg, and concludes as an adult that reproduces, producing an offspring in the form of a new zygote which then itself goes through the same series of stages, the process repeating in a cyclic fashion.
"The concept is closely related to those of the life history, development and ontogeny, but differs from them in stressing renewal." Transitions of form may involve growth, asexual reproduction, or sexual reproduction.
In some organisms, different "generations" of the species succeed each other during the life cycle. For plants and many algae, there are two multicellular stages, and the life cycle is referred to as alternation of generations. The term life history is often used, particularly for organisms such as the red algae which have three multicellular stages (or more), rather than two.
Life cycles that include sexual reproduction involve alternating haploid (n) and diploid (2n) stages, i.e., a change of ploidy is involved. To return from a diploid stage to a haploid stage, meiosis must occur. In regard to changes of ploidy, there are three types of cycles:
haplontic life cycle — the haploid stage is multicellular and the diploid stage is a single cell, meiosis is "zygotic".
diplontic life cycle — the diploid stage is multicellular and haploid gametes are formed, meiosis is "gametic".
haplodiplontic life cycle (also referred to as diplohaplontic, diplobiontic, or dibiontic life cycle) — multicellular diploid and haploid stages occur, meiosis is "sporic".
The cycles differ in when mitosis (growth) occurs. Zygotic meiosis and gametic meiosis have one mitotic stage: mitosis occurs during the n phase in zygotic meiosis and during the 2n phase in gametic meiosis. Therefore, zygotic and gametic meiosis are collectively termed "haplobiontic" (single mitotic phase, not to be confused with haplontic). Sporic meiosis, on the other hand, has mitosis in two stages, both the diploid and haploid stages, termed "diplobiontic" (not to be confused with diplontic).
Discovery
The study of reproduction and development in organisms was carried out by many botanists and zoologists.
Wilhelm Hofmeister demonstrated that alternation of generations is a feature that unites plants, and published this result in 1851 (see plant sexuality).
Some terms (haplobiont and diplobiont) used for the description of life cycles were proposed initially for algae by Nils Svedelius, and then became used for other organisms. Other terms (autogamy and gamontogamy) used in protist life cycles were introduced by Karl Gottlieb Grell. The description of the complex life cycles of various organisms contributed to the disproof of the ideas of spontaneous generation in the 1840s and 1850s.
Haplontic life cycle
A zygotic meiosis is a meiosis of a zygote immediately after karyogamy, which is the fusion of two cell nuclei. This way, the organism ends its diploid phase and produces several haploid cells. These cells divide mitotically to form either larger, multicellular individuals, or more haploid cells. Two opposite types of gametes (e.g., male and female) from these individuals or cells fuse to become a zygote.
In the whole cycle, zygotes are the only diploid cell; mitosis occurs only in the haploid phase.
The individuals or cells as a result of mitosis are haplonts, hence this life cycle is also called haplontic life cycle. Haplonts are:
In archaeplastidans: some green algae (e.g., Chlamydomonas, Zygnema, Chara)
In stramenopiles: some golden algae
In alveolates: many dinoflagellates, e.g., Ceratium, Gymnodinium, some apicomplexans (e.g., Plasmodium)
In rhizarians: some euglyphids, ascetosporeans
In excavates: some parabasalids
In amoebozoans: Dictyostelium
In opisthokonts: most fungi (some chytrids, zygomycetes, some ascomycetes, basidiomycetes)
Diplontic life cycle
In gametic meiosis, instead of immediately dividing meiotically to produce haploid cells, the zygote divides mitotically to produce a multicellular diploid individual or a group of more unicellular diploid cells. Cells from the diploid individuals then undergo meiosis to produce haploid cells or gametes. Haploid cells may divide again (by mitosis) to form more haploid cells, as in many yeasts, but the haploid phase is not the predominant life cycle phase. In most diplonts, mitosis occurs only in the diploid phase, i.e. gametes usually form quickly and fuse to produce diploid zygotes.
In the whole cycle, gametes are usually the only haploid cells, and mitosis usually occurs only in the diploid phase.
The diploid multicellular individual is a diplont, hence a gametic meiosis is also called a diplontic life cycle. Diplonts are:
In archaeplastidans: some green algae (e.g., Cladophora glomerata, Acetabularia)
In stramenopiles: some brown algae (the Fucales, however, their life cycle can also be interpreted as strongly heteromorphic-diplohaplontic, with a highly reduced gametophyte phase, as in the flowering plants), some xanthophytes (e.g., Vaucheria), most diatoms, some oomycetes (e.g., Saprolegnia, Plasmopara viticola), opalines, some "heliozoans" (e.g., Actinophrys, Actinosphaerium)
In alveolates: ciliates
In excavates: some parabasalids
In opisthokonts: animals, some fungi (e.g., some ascomycetes)
Haplodiplontic life cycle
In sporic meiosis (also commonly known as intermediary meiosis), the zygote divides mitotically to produce a multicellular diploid sporophyte. The sporophyte creates spores via meiosis which also then divide mitotically producing haploid individuals called gametophytes. The gametophytes produce gametes via mitosis. In some plants the gametophyte is not only small-sized but also short-lived; in other plants and many algae, the gametophyte is the "dominant" stage of the life cycle.
Haplodiplonts are:
In archaeplastidans: red algae (which have two sporophyte generations), some green algae (e.g., Ulva), land plants
In stramenopiles: most brown algae
In rhizarians: many foraminiferans, plasmodiophoromycetes
In amoebozoa: myxogastrids
In opisthokonts: some fungi (some chytrids, some ascomycetes like the brewer's yeast)
Other eukaryotes: haptophytes
Some animals have a sex-determination system called haplodiploid, but this is not related to the haplodiplontic life cycle.
Vegetative meiosis
Some red algae (such as Bonnemaisonia and Lemanea) and green algae (such as Prasiola) have vegetative meiosis, also called somatic meiosis, which is a rare phenomenon. Vegetative meiosis can occur in haplodiplontic and also in diplontic life cycles. The gametophytes remain attached to and part of the sporophyte. Vegetative (non-reproductive) diploid cells undergo meiosis, generating vegetative haploid cells. These undergo many mitosis, and produces gametes.
A different phenomenon, called vegetative diploidization, a type of apomixis, occurs in some brown algae (e.g., Elachista stellaris). Cells in a haploid part of the plant spontaneously duplicate their chromosomes to produce diploid tissue.
Parasitic life cycle
Parasites depend on the exploitation of one or more hosts. Those that must infect more than one host species to complete their life cycles are said to have complex or indirect life cycles. Dirofilaria immitis, or the heartworm, has an indirect life cycle, for example. The microfilariae must first be ingested by a female mosquito, where it develops into the infective larval stage. The mosquito then bites an animal and transmits the infective larvae into the animal, where they migrate to the pulmonary artery and mature into adults.
Those parasites that infect a single species have direct life cycles. An example of a parasite with a direct life cycle is Ancylostoma caninum, or the canine hookworm. They develop to the infective larval stage in the environment, then penetrate the skin of the dog directly and mature to adults in the small intestine.
If a parasite has to infect a given host in order to complete its life cycle, then it is said to be an obligate parasite of that host; sometimes, infection is facultative—the parasite can survive and complete its life cycle without infecting that particular host species. Parasites sometimes infect hosts in which they cannot complete their life cycles; these are accidental hosts.
A host in which parasites reproduce sexually is known as the definitive, final or primary host. In intermediate hosts, parasites either do not reproduce or do so asexually, but the parasite always develops to a new stage in this type of host. In some cases a parasite will infect a host, but not undergo any development, these hosts are known as paratenic or transport hosts. The paratenic host can be useful in raising the chance that the parasite will be transmitted to the definitive host. For example, the cat lungworm (Aelurostrongylus abstrusus) uses a slug or snail as an intermediate host; the first stage larva enters the mollusk and develops to the third stage larva, which is infectious to the definitive host—the cat. If a mouse eats the slug, the third stage larva will enter the mouse's tissues, but will not undergo any development.
Evolution
The primitive type of life cycle probably had haploid individuals with asexual reproduction. Bacteria and archaea exhibit a life cycle like this, and some eukaryotes apparently do too (e.g., Cryptophyta, Choanoflagellata, many Euglenozoa, many Amoebozoa, some red algae, some green algae, the imperfect fungi, some rotifers and many other groups, not necessarily haploid). However, these eukaryotes probably are not primitively asexual, but have lost their sexual reproduction, or it just was not observed yet. Many eukaryotes (including animals and plants) exhibit asexual reproduction, which may be facultative or obligate in the life cycle, with sexual reproduction occurring more or less frequently.
Individual organisms participating in a biological life cycle ordinarily age and die, while cells from these organisms that connect successive life cycle generations (germ line cells and their descendants) are potentially immortal. The basis for this difference is a fundamental problem in biology. The Russian biologist and historian Zhores A. Medvedev considered that the accuracy of genome replicative and other synthetic systems alone cannot explain the immortality of germlines. Rather Medvedev thought that known features of the biochemistry and genetics of sexual reproduction indicate the presence of unique information maintenance and restoration processes at the gametogenesis stage of the biological life cycle. In particular, Medvedev considered that the most important opportunities for information maintenance of germ cells are created by recombination during meiosis and DNA repair; he saw these as processes within the germ line cells that were capable of restoring the integrity of DNA and chromosomes from the types of damage that cause irreversible ageing in non-germ line cells, e.g. somatic cells.
The ancestry of each present day cell presumably traces back, in an unbroken lineage for over 3 billion years to the origin of life. It is not actually cells that are immortal but multi-generational cell lineages. The immortality of a cell lineage depends on the maintenance of cell division potential. This potential may be lost in any particular lineage because of cell damage, terminal differentiation as occurs in nerve cells, or programmed cell death (apoptosis) during development. Maintenance of cell division potential of the biological life cycle over successive generations depends on the avoidance and the accurate repair of cellular damage, particularly DNA damage. In sexual organisms, continuity of the germline over successive cell cycle generations depends on the effectiveness of processes for avoiding DNA damage and repairing those DNA damages that do occur. Sexual processes in eukaryotes provide an opportunity for effective repair of DNA damages in the germ line by homologous recombination.
| Biology and health sciences | Ecology | Biology |
319653 | https://en.wikipedia.org/wiki/Bidet | Bidet | A bidet ( or ) is a bowl or receptacle designed to be sat upon in order to wash a person's genitalia, perineum, inner buttocks, and anus. The modern variety has a plumbed-in water supply and a drainage opening, and is thus a plumbing fixture subject to local hygiene regulations. The bidet is designed to promote personal hygiene and is used after defecation, and before and after sexual intercourse. It can also be used to wash feet, with or without filling it up with water. Some people even use bidets to bathe babies or pets. In several European countries, a bidet is now required by law to be present in every bathroom containing a toilet bowl. It was originally located in the bedroom, near the chamber-pot and the marital bed, but in modern times is located near the toilet bowl in the bathroom. Fixtures that combine a toilet seat with a washing facility include the electronic bidet.
Opinions as to the necessity of the bidet vary widely over different nationalities and cultures. In cultures that use it habitually, such as parts of Western, Central and Southeastern Europe (especially Italy and Portugal), Eastern Asia and some Latin American countries such as Argentina or Paraguay, it is considered an indispensable tool in maintaining good personal hygiene. It is commonly used in North African countries, such as Egypt. It is rarely used in sub-Saharan Africa, Australia, and North America.
"Bidet" is a French loanword meaning "pony" due to the straddling position adopted in its usage.
Applications
Bidets are primarily used to wash and clean the genitalia, perineum, inner buttocks, and anus. Some bidets have a vertical jet intended to give easy access for washing and rinsing the perineum and anal area. The traditional separate bidet is like a wash-basin which is used with running warm water with the help of specific soaps, and may then be used for many other purposes such as washing feet.
Types
Bidet shower
A bidet shower (also known as "bidet spray", "bidet sprayer", or "health faucet") is a hand-held triggered nozzle, similar to that on a kitchen sink sprayer, that delivers a spray of water to assist in anal cleansing and cleaning the genitals after defecation and urination. In contrast to a bidet that is integrated with the toilet, a bidet shower has to be held by the hands, and cleaning does not take place automatically. Bidet showers are common in countries where water is considered essential for anal cleansing.
Drawbacks include the possibility of wetting a user's clothing if used carelessly. In addition, a user must be reasonably mobile and flexible to use a hand-held bidet shower.
Conventional or standalone bidet
A bidet is a plumbing fixture that is installed as a separate unit in the bathroom besides toilet, shower and sink, which users have to straddle. Some bidets resemble a large hand basin, with taps and a stopper so they can be filled up; other designs have a nozzle that squirts a jet of water to aid in cleansing.
Add-on bidets
There are bidets that are attachable to toilet bowls, saving space and obviating additional plumbing. A bidet may be a movable or fixed nozzle, either attached to an existing toilet on the back or side toilet rim, or replacing the toilet seat. In these cases, its use is restricted to cleaning the anus and genitals. Some bidets of this type produce a vertical water jet and others a more-or-less oblique one. Other bidets have one nozzle on the side rim aimed at both anal and genital areas, and other designs have two nozzles on the back rim. The shorter one, called the "family nozzle", is used for washing the area around the anus, and the longer one ("bidet nozzle") is designed for washing the vulva.
Such attachable bidets (also called "combined toilets", "bidet attachments", or "add-on bidets") are controlled either mechanically, by turning a valve, or electronically. Electronic bidets are controlled with waterproof electrical switches rather than a manual valve. There are models that have a heating element which blows warm air to dry the user after washing, that offer heated seats, wireless remote controls, illumination through built in night lights, or built in deodorizers and activated carbon filters to remove odours. Further refinements include adjustable water pressure, temperature compensation, and directional spray control. Where bathroom appearance is of concern, under-the-seat mounting types have become more popular.
An add-on bidet typically connects to the existing water supply of a toilet via the addition of a threaded tee pipe adapter, and requires no soldering or other plumbing work. Electronic add-on bidets also require a GFCI protected grounded electrical outlet.
Usage and health
Personal hygiene is improved and maintained more accurately and easily with the use of both toilet paper and a bidet as compared to the use of toilet paper alone. In some add-on bidets with vertical jets, little water is used and toilet paper may not be necessary. Addressing hemorrhoids and genital health issues might also be facilitated by the use of bidet fixtures.
Because of the large surface of the basin, after-use and routine disinfection of stand-alone bidets require thoroughness, or microbial contamination from one user to the next could take place. Bidet attachments are sometimes included on hospital toilets because of their utility in maintaining hygiene. Hospitals must consider the use of bidet properly and consider the clinical background of patients to prevent cross-infection. Warm-water bidets may harbor dangerous microbes if not properly disinfected.
Environmental aspects
From an environmental standpoint, bidets can reduce the need for toilet paper. Considering that an average person uses only 0.5 litre (1/8 US gallon) of water for cleansing by using a bidet, much less water is used than for manufacturing toilet paper. An article in Scientific American concluded that using a bidet is "much less stressful on the environment than using paper". Scientific American has also reported that if the US switched to using bidets, 15 million trees could be saved every year.
In the US, UK, and some other countries, wet wipes are heavily marketed as an upgrade from dry toilet paper. However, this product has been criticized for its adverse environmental impact, due to the non-biodegradable plastic fibers composing most versions. Although the wipes are promoted as "flushable", they absorb waste fats and agglomerate into massive "fatbergs" which can clog sewer systems and must be cleared at great expense. Bidets are being marketed as cleaning better than toilet paper or wet wipes, with fewer negative environmental effects.
Society and culture
The bidet is common in Catholic countries and required by law in some. It is also found in some traditionally Eastern Orthodox and Protestant countries such as Greece and Finland respectively, where bidet showers are common.
In Islam, there are many strict rules concerning excretion; in particular, anal washing with water is required. Consequently, in Middle Eastern regions where Islam is the predominant religion, water for anal washing is provided in most toilets, usually in the form of a hand-held "bidet shower" or shattaf.
Prevalence
Bidets are becoming increasingly popular with the elderly and disabled. Combined toilet/bidet installations make self-care toileting possible for many people, affording greater independence. There are often special units with higher toilet seats allowing easier wheelchair transfer, and with some form of electronic remote control that benefits an individual with limited mobility or otherwise requiring assistance.
Bidets are common bathroom fixtures in the Arab world and in Catholic countries, such as Italy (the installation of a bidet in a bathroom has been mandatory since 1975), Spain (but in recent times new or renewed houses tend to have bathrooms without bidets, except the luxurious ones), and Portugal (installation is mandatory since 1975). They are also found in Southeastern European countries such as Albania, Bosnia and Herzegovina, Romania, Greece and Turkey. They are very popular in some South American countries, particularly Argentina, Paraguay and Uruguay. Electronic bidet-integrated toilets, often with functions such as toilet seat warming, are commonly found in Japan, and are becoming more popular in other Asian countries.
In Northern Europe, bidets are rare, although in Finland, bidet showers are common. Bidet showers are most commonly found in Southeast Asia, South Asia, and the Middle East.
In 1980, the first "paperless toilet" was launched in Japan by manufacturer Toto, a combination of toilet and bidet which also dries the user after washing. These combination toilet-bidets (washlet) with seat warmers, or attachable bidets are particularly popular in Japan and South Korea, and are found in approximately 76% of Japanese households . They are commonly found in hotels and some public facilities. These bidet-toilets, along with toilet seat and bidet units (to convert an existing toilet) are sold in many countries, including the United States.
Bidet seat conversions are much easier and lower cost to install than traditional bidets, and have disrupted the market for the older fixtures.
After a slow start in the 1990s, electronic bidets are starting to become more available in the United States. American distributors were directly influenced by their Japanese predecessors, as the founders of Brondell (established in 2003) have indicated. The popularity of add-on bidet units is steadily increasing in the United States, Canada and the United Kingdom, in part because of their ability to treat hemorrhoids or urogenital infections. In addition, shortages of toilet paper due to the coronavirus pandemic have led to an increased interest in bidets.
Etymology
Bidet is a French word for "pony", and in Old French, meant "to trot". This etymology comes from the notion that one "rides" or straddles a bidet much like a pony is ridden. The word "bidet" was used in 15th-century France to refer to the pet ponies that French royalty kept.
History
The bidet appears to have been an invention of French furniture makers in the late 17th century, although no exact date or inventor is known. The earliest written reference to the bidet is in 1726 in Italy. Even though there are records of Maria Carolina of Austria, Queen of Naples and Sicily, requesting a bidet for her personal bathroom in the Royal Palace of Caserta in the second half of the 18th century, the bidet did not become widespread in Italy until after the Second World War. The bidet is possibly associated with the chamber pot and the bourdaloue, the latter being a small, hand-held chamber pot.
Historical antecedents and early functions of the bidet are believed to include devices used for contraception. Bidets are considered ineffective by today's standards of contraception, and their use for that function was quickly abandoned and forgotten following the advent of modern contraceptives such as the pill.
By 1900, due to plumbing improvements, the bidet (and chamber pot) moved from the bedroom to the bathroom and became more convenient to fill and drain.
In 1928, in the United States, John Harvey Kellogg applied for a patent on an "anal douche". In his application, he used the term to describe a system comparable to what today might be called a bidet nozzle, which can be attached to a toilet to perform anal cleansing with water.
In 1965, the American Bidet Company featured an adjustable spray nozzle and warm water option, seeking to make the bidet a household item. The fixture was expensive, and required floor space to install; it was eventually discontinued without a replacement model.
The early 1980s saw the introduction of the electronic bidet from Japan, with names such as Clean Sense, Galaxy, Infinity, Novita, and of non-electric attachments such as Gobidet. These devices have attachments that connect to existing toilet water supplies, and can be used in bathrooms lacking the space for a separate bidet and toilet. Many models have additional features, such as instant-heating warm water, night lights, or a heated seat.
| Technology | Hydraulics and pneumatics | null |
319712 | https://en.wikipedia.org/wiki/Well-founded%20relation | Well-founded relation | In mathematics, a binary relation is called well-founded (or wellfounded or foundational) on a set or, more generally, a class if every non-empty subset has a minimal element with respect to ; that is, there exists an such that, for every , one does not have . In other words, a relation is well-founded if:
Some authors include an extra condition that is set-like, i.e., that the elements less than any given element form a set.
Equivalently, assuming the axiom of dependent choice, a relation is well-founded when it contains no infinite descending chains, which can be proved when there is no infinite sequence of elements of such that for every natural number .
In order theory, a partial order is called well-founded if the corresponding strict order is a well-founded relation. If the order is a total order then it is called a well-order.
In set theory, a set is called a well-founded set if the set membership relation is well-founded on the transitive closure of . The axiom of regularity, which is one of the axioms of Zermelo–Fraenkel set theory, asserts that all sets are well-founded.
A relation is converse well-founded, upwards well-founded or Noetherian on , if the converse relation is well-founded on . In this case is also said to satisfy the ascending chain condition. In the context of rewriting systems, a Noetherian relation is also called terminating.
Induction and recursion
An important reason that well-founded relations are interesting is because a version of transfinite induction can be used on them: if () is a well-founded relation, is some property of elements of , and we want to show that
holds for all elements of ,
it suffices to show that:
If is an element of and is true for all such that , then must also be true.
That is,
Well-founded induction is sometimes called Noetherian induction, after Emmy Noether.
On par with induction, well-founded relations also support construction of objects by transfinite recursion. Let be a set-like well-founded relation and a function that assigns an object to each pair of an element and a function on the initial segment of . Then there is a unique function such that for every ,
That is, if we want to construct a function on , we may define using the values of for .
As an example, consider the well-founded relation , where is the set of all natural numbers, and is the graph of the successor function . Then induction on is the usual mathematical induction, and recursion on gives primitive recursion. If we consider the order relation , we obtain complete induction, and course-of-values recursion. The statement that is well-founded is also known as the well-ordering principle.
There are other interesting special cases of well-founded induction. When the well-founded relation is the usual ordering on the class of all ordinal numbers, the technique is called transfinite induction. When the well-founded set is a set of recursively-defined data structures, the technique is called structural induction. When the well-founded relation is set membership on the universal class, the technique is known as ∈-induction. See those articles for more details.
Examples
Well-founded relations that are not totally ordered include:
The positive integers , with the order defined by if and only if divides and .
The set of all finite strings over a fixed alphabet, with the order defined by if and only if is a proper substring of .
The set of pairs of natural numbers, ordered by if and only if and .
Every class whose elements are sets, with the relation ∈ ("is an element of"). This is the axiom of regularity.
The nodes of any finite directed acyclic graph, with the relation defined such that if and only if there is an edge from to .
Examples of relations that are not well-founded include:
The negative integers , with the usual order, since any unbounded subset has no least element.
The set of strings over a finite alphabet with more than one element, under the usual (lexicographic) order, since the sequence is an infinite descending chain. This relation fails to be well-founded even though the entire set has a minimum element, namely the empty string.
The set of non-negative rational numbers (or reals) under the standard ordering, since, for example, the subset of positive rationals (or reals) lacks a minimum.
Other properties
If is a well-founded relation and is an element of , then the descending chains starting at are all finite, but this does not mean that their lengths are necessarily bounded. Consider the following example:
Let be the union of the positive integers with a new element ω that is bigger than any integer. Then is a well-founded set, but
there are descending chains starting at ω of arbitrary great (finite) length;
the chain has length for any .
The Mostowski collapse lemma implies that set membership is a universal among the extensional well-founded relations: for any set-like well-founded relation on a class that is extensional, there exists a class such that is isomorphic to .
Reflexivity
A relation is said to be reflexive if holds for every in the domain of the relation. Every reflexive relation on a nonempty domain has infinite descending chains, because any constant sequence is a descending chain. For example, in the natural numbers with their usual order ≤, we have . To avoid these trivial descending sequences, when working with a partial order ≤, it is common to apply the definition of well foundedness (perhaps implicitly) to the alternate relation < defined such that if and only if and . More generally, when working with a preorder ≤, it is common to use the relation < defined such that if and only if and . In the context of the natural numbers, this means that the relation <, which is well-founded, is used instead of the relation ≤, which is not. In some texts, the definition of a well-founded relation is changed from the definition above to include these conventions.
| Mathematics | Order theory | null |
319834 | https://en.wikipedia.org/wiki/Rutherford%20scattering%20experiments | Rutherford scattering experiments | The Rutherford scattering experiments were a landmark series of experiments by which scientists learned that every atom has a nucleus where all of its positive charge and most of its mass is concentrated. They deduced this after measuring how an alpha particle beam is scattered when it strikes a thin metal foil. The experiments were performed between 1906 and 1913 by Hans Geiger and Ernest Marsden under the direction of Ernest Rutherford at the Physical Laboratories of the University of Manchester.
The physical phenomenon was explained by Rutherford in a classic 1911 paper that eventually lead to the widespread use of scattering in particle physics to study subatomic matter. Rutherford scattering or Coulomb scattering is the elastic scattering of charged particles by the Coulomb interaction. The paper also initiated the development of the planetary Rutherford model of the atom and eventually the Bohr model.
Rutherford scattering is now exploited by the materials science community in an analytical technique called Rutherford backscattering.
Summary
Thomson's model of the atom
The prevailing model of atomic structure before Rutherford's experiments was devised by J. J. Thomson. Thomson had discovered the electron through his work on cathode rays and proposed that they existed within atoms, and an electric current is electrons hopping from one atom to an adjacent one in a series. There logically had to be a commensurate amount of positive charge to balance the negative charge of the electrons and hold those electrons together. Having no idea what the source of this positive charge was, he tentatively proposed that the positive charge was everywhere in the atom, adopting a spherical shape for simplicity. Thomson imagined that the balance of electrostatic forces would distribute the electrons throughout this sphere in a more or less even manner. Thomson also believed the electrons could move around in this sphere, and in that regard he likened the substance of the sphere to a liquid. In fact the positive sphere was more of an abstraction than anything material. He did not propose a positively-charged subatomic particle; a counterpart to the electron.
Thomson was never able to develop a complete and stable model that could predict any of the other known properties of the atom, such as emission spectra and valencies. The Japanese scientist Hantaro Nagaoka rejected Thomson's model on the grounds that opposing charges cannot penetrate each other. He proposed instead that electrons orbit the positive charge like the rings around Saturn. However this model was also known to be unstable.
Alpha particles and the Thomson atom
An alpha particle is a positively charged particle of matter that is spontaneously emitted from certain radioactive elements. Alpha particles are so tiny as to be invisible, but they can be detected with the use of phosphorescent screens, photographic plates, or electrodes. Rutherford discovered them in 1899. In 1906, by studying how alpha particle beams are deflected by magnetic and electric fields, he deduced that they were essentially helium atoms stripped of two electrons. Thomson and Rutherford knew nothing about the internal structure of alpha particles. Prior to 1911 they were thought to have a diameter similar to helium atoms and contain ten or so electrons.
Thomson's model was consistent with the experimental evidence available at the time. Thomson studied beta particle scattering which showed small angle deflections modelled as interactions of the particle with many atoms in succession. Each interaction of the particle with the electrons of the atom and the positive background sphere would lead to a tiny deflection, but many such collisions could add up. The scattering of alpha particles was expected to be similar. Rutherford's team would show that the multiple scattering model was not needed: single scattering from a compact charge at the centre of the atom would account for all of the scattering data.
Rutherford, Geiger, and Marsden
Ernest Rutherford was Langworthy Professor of Physics at the Victoria University of Manchester (now the University of Manchester). He had already received numerous honours for his studies of radiation. He had discovered the existence of alpha rays, beta rays, and gamma rays, and had proved that these were the consequence of the disintegration of atoms. In 1906, he received a visit from the German physicist Hans Geiger, and was so impressed that he asked Geiger to stay and help him with his research. Ernest Marsden was a physics undergraduate student studying under Geiger.
In 1908, Rutherford sought to independently determine the charge and mass of alpha particles. To do this, he wanted to count the number of alpha particles and measure their total charge; the ratio would give the charge of a single alpha particle. Alpha particles are too tiny to see, but Rutherford knew about Townsend discharge, a cascade effect from ionisation leading to a pulse of electric current. On this principle, Rutherford and Geiger designed a simple counting device which consisted of two electrodes in a glass tube. (See #1908 experiment.) Every alpha particle that passed through the tube would create a pulse of electricity that could be counted. It was an early version of the Geiger counter.
The counter that Geiger and Rutherford built proved unreliable because the alpha particles were being too strongly deflected by their collisions with the molecules of air within the detection chamber. The highly variable trajectories of the alpha particles meant that they did not all generate the same number of ions as they passed through the gas, thus producing erratic readings. This puzzled Rutherford because he had thought that alpha particles were too heavy to be deflected so strongly. Rutherford asked Geiger to investigate how far matter could scatter alpha rays.
The experiments they designed involved bombarding a metal foil with a beam of alpha particles to observe how the foil scattered them in relation to its thickness and material. They used a phosphorescent screen to measure the trajectories of the particles. Each impact of an alpha particle on the screen produced a tiny flash of light. Geiger worked in a darkened lab for hours on end, counting these tiny scintillations using a microscope. For the metal foil, they tested a variety of metals, but favoured gold because they could make the foil very thin, as gold is the most malleable metal. As a source of alpha particles, Rutherford's substance of choice was radium, which is thousands of times more radioactive than uranium.
Scattering theory and the new atomic model
In a 1909 experiment, Geiger and Marsden discovered that the metal foils could scatter some alpha particles in all directions, sometimes more than 90°. This should have been impossible according to Thomson's model. According to Thomson's model, all the alpha particles should have gone straight through.
In Thomson's model of the atom, the sphere of positive charge that fills the atom and encapsulates the electrons is permeable; the electrons could move around in it, after all. Therefore, an alpha particle should be able to pass through this sphere if the electrostatic forces within permit it. Thomson himself did not study how an alpha particle might be scattered in such a collision with an atom, but he did study beta particle scattering. He calculated that a beta particle would only experience very small deflection when passing through an atom, and even after passing through many atoms in a row, the total deflection should still be less than 1°. Alpha particles typically have much more momentum than beta particles and therefore should likewise experience only the slightest deflection.
The extreme scattering observed forced Rutherford to revise the model of the atom. The issue in Thomson's model was that the charges were too diffuse to produce a sufficiently strong electrostatic force to cause such repulsion. Therefore they had to be more concentrated. In Rutherford's new model, the positive charge does not fill the entire volume of the atom but instead constitutes a tiny nucleus at least 10,000 times smaller than the atom as a whole. All that positive charge concentrated in a much smaller volume produces a much stronger electric field near its surface. The nucleus also carried most of the atom's mass. This meant that it could deflect alpha particles by up to 180° depending on how close they pass. The electrons surround this nucleus, spread throughout the atom's volume. Because their negative charge is diffuse and their combined mass is low, they have a negligible effect on the alpha particle.
To verify his model, Rutherford developed a scientific model to predict the intensity of alpha particles at the different angles they scattered coming out of the gold foil, assuming all of the positive charge was concentrated at the centre of the atom. This model was validated in an experiment performed in 1913. His model explained both the beta scattering results of Thomson and the alpha scattering results of Geiger and Marsden.
Legacy
There was little reaction to Rutherford's now-famous 1911 paper in the first years. The paper was primarily about alpha particle scattering in an era before particle scattering was a primary tool for physics. The probability techniques he used and confusing collection of observations involved were not immediately compelling.
Nuclear physics
The first impacts were to encourage new focus on scattering experiments. For example the first results from a cloud chamber, by C.T.R. Wilson shows alpha particle scattering and also appeared in 1911. Over time, particle scattering became a major aspect of theoretical and experimental physics; Rutherford's concept of a "cross-section" now dominates the descriptions of experimental particle physics. The historian Silvan S. Schweber suggests that Rutherford's approach marked the shift to viewing all interactions and measurements in physics as scattering processes. After the nucleus - a term Rutherford introduced in 1912 - became the accepted model for the core of atoms, Rutherford's analysis of the scattering of alpha particles created a new branch of physics, nuclear physics.
Atomic model
Rutherford's new atom model caused no stir. Rutherford explicitly ignores the electrons, only mentioning Hantaro Nagaoka's Saturnian model of electrons orbiting a tiny "sun", a model that had been previously rejected as mechanically unstable. By ignoring the electrons Rutherford also ignores any potential implications for atomic spectroscopy for chemistry. Rutherford himself did not press the case for his atomic model: his own 1913 book on "Radioactive substances and their radiations" only mentions the atom twice; other books by other authors around this time focus on Thomson's model.
The impact of Rutherford's nuclear model came after Niels Bohr arrived as a post-doctoral student in Manchester at Rutherford's invitation. Bohr dropped his work on the Thomson model in favour of Rutherford's nuclear model, developing the Rutherford–Bohr model over the next several years. Eventually Bohr incorporated early ideas of quantum mechanics into the model of the atom, allowing prediction of electronic spectra and concepts of chemistry.
Hantaro Nagaoka, who had proposed a Saturnian model of the atom, wrote to Rutherford from Tokyo in 1911: "I have been struck with the simpleness of the apparatus you employ and the brilliant results you obtain." The astronomer Arthur Eddington called Rutherford's discovery the most important scientific achievement since Democritus proposed the atom ages earlier. Rutherford has since been hailed as "the father of nuclear physics".
In a lecture delivered on 15 October 1936 at Cambridge University, Rutherford described his shock at the results of the 1909 experiment:
Rutherford's claim of surprise makes a good story but by the time of the Geiger-Marsden experiment the result confirmed suspicions Rutherford developed from his many previous experiments.
Experiments
Alpha particle scattering: 1906 and 1908 experiments
Rutherford's first steps towards his discovery of the nature of the atom came from his work to understand alpha particles.
In 1906, Rutherford noticed that alpha particles passing through sheets of mica were deflected by the sheets by as much as 2 degrees. Rutherford placed a radioactive source in a sealed tube ending with a narrow slits followed by a photographic plate. Half of the slit was covered by a thin layer of mica. A magnetic field around the tube was altered every 10 minutes to reject the effect of beta rays, known to be sensitive to magnetic fields. The tube was evacuated to different amounts and a series of images recorded. At the lowest pressure the image of the open slit was clear, while images of the mica covered slit or the open slit at higher pressures were fuzzy. Rutherford explained these results as alpha-particle scattering in a paper published in 1906. He already understood the implications of the observation for models of atoms: "such a result brings out clearly the fact that the atoms of matter must be the seat of very intense electrical forces".
A 1908 paper by Geiger, On the Scattering of α-Particles by Matter, describes the following experiment. He constructed a long glass tube, nearly two metres long. At one end of the tube was a quantity of "radium emanation" (R) as a source of alpha particles. The opposite end of the tube was covered with a phosphorescent screen (Z). In the middle of the tube was a 0.9 mm-wide slit. The alpha particles from R passed through the slit and created a glowing patch of light on the screen. A microscope (M) was used to count the scintillations on the screen and measure their spread. Geiger pumped all the air out of the tube so that the alpha particles would be unobstructed, and they left a neat and tight image on the screen that corresponded to the shape of the slit. Geiger then allowed some air into the tube, and the glowing patch became more diffuse. Geiger then pumped out the air and placed one or two gold foils over the slit at AA. This too caused the patch of light on the screen to become more spread out, with the larger spread for two layers. This experiment demonstrated that both air and solid matter could markedly scatter alpha particles.
Alpha particle reflection: the 1909 experiment
The results of the initial alpha particle scattering experiments were confusing. The angular spread of the particle on the screen varied greatly with the shape of the apparatus and its internal pressure. Rutherford suggested that Ernest Marsden, a physics undergraduate student studying under Geiger, should look for diffusely reflected or back-scattered alpha particles, even though these were not expected. Marsden's first crude reflector got results, so Geiger helped him create a more sophisticated apparatus. They were able to demonstrate that 1 in 8000 alpha particle collisions were diffuse reflections. Although this fraction was small, it was much larger than the Thomson model of the atom could explain.
These results where published in a 1909 paper, On a Diffuse Reflection of the α-Particles, where Geiger and Marsden described the experiment by which they proved that alpha particles can indeed be scattered by more than 90°. In their experiment, they prepared a small conical glass tube (AB) containing "radium emanation" (radon), "radium A" (actual radium), and "radium C" (bismuth-214); its open end was sealed with mica. This was their alpha particle emitter. They then set up a lead plate (P), behind which they placed a fluorescent screen (S). The tube was held on the opposite side of plate, such that the alpha particles it emitted could not directly strike the screen. They noticed a few scintillations on the screen because some alpha particles got around the plate by bouncing off air molecules. They then placed a metal foil (R) to the side of the lead plate. They tested with lead, gold, tin, aluminium, copper, silver, iron, and platinum. They pointed the tube at the foil to see if the alpha particles would bounce off it and strike the screen on the other side of the plate, and observed an increase in the number of scintillations on the screen. Counting the scintillations, they observed that metals with higher atomic mass, such as gold, reflected more alpha particles than lighter ones such as aluminium.
Geiger and Marsden then wanted to estimate the total number of alpha particles that were reflected. The previous setup was unsuitable for doing this because the tube contained several radioactive substances (radium plus its decay products) and thus the alpha particles emitted had varying ranges, and because it was difficult for them to ascertain at what rate the tube was emitting alpha particles. This time, they placed a small quantity of radium C (bismuth-214) on the lead plate, which bounced off a platinum reflector (R) and onto the screen. They concluded that approximately 1 in 8,000 of the alpha particles that struck the reflector bounced onto the screen. By measuring the reflection from thin foils they showed that the effect due to a volume and not a surface effect. When contrasted with the vast number of alpha particles that pass unhindered through a metal foil, this small number of large angle reflections was a strange result that meant very large forces were involved.
Dependence on foil material and thickness: the 1910 experiment
A 1910 paper by Geiger, The Scattering of the α-Particles by Matter, describes an experiment to measure how the most probable angle through which an alpha particle is deflected varies with the material it passes through, the thickness of the material, and the velocity of the alpha particles. He constructed an airtight glass tube from which the air was pumped out. At one end was a bulb (B) containing "radium emanation" (radon-222). By means of mercury, the radon in B was pumped up the narrow glass pipe whose end at A was plugged with mica. At the other end of the tube was a fluorescent zinc sulfide screen (S). The microscope which he used to count the scintillations on the screen was affixed to a vertical millimetre scale with a vernier, which allowed Geiger to precisely measure where the flashes of light appeared on the screen and thus calculate the particles' angles of deflection. The alpha particles emitted from A was narrowed to a beam by a small circular hole at D. Geiger placed a metal foil in the path of the rays at D and E to observe how the zone of flashes changed. He tested gold, tin, silver, copper, and aluminium. He could also vary the velocity of the alpha particles by placing extra sheets of mica or aluminium at A.
From the measurements he took, Geiger came to the following conclusions:
the most probable angle of deflection increases with the thickness of the material
the most probable angle of deflection is proportional to the atomic mass of the substance
the most probable angle of deflection decreases with the velocity of the alpha particles
Rutherford's Structure of the Atom paper (1911)
Considering the results of these experiments, Rutherford published a landmark paper in 1911 titled "The Scattering of α and β Particles by Matter and the Structure of the Atom" wherein he showed that single scattering from a very small and intense electric charge predicts primarily small-angle scattering with small but measurable amounts of backscattering. For the purpose of his mathematical calculations he assumed this central charge was positive, but he admitted he could not prove this and that he had to wait for other experiments to develop his theory.
Rutherford developed a mathematical equation that modelled how the foil should scatter the alpha particles if all the positive charge and most of the atomic mass was concentrated in a point at the centre of an atom.
From the scattering data, Rutherford estimated the central charge qn to be about +100 units.
Rutherford's paper does not discuss any electron arrangement beyond discussions on the scattering from Thomson's plum pudding model and Nagaoka's Saturnian model. He shows that the scattering results predicted by Thomson's model are also explained by single scattering, but that Thomson's model does not explain large angle scattering. He says that Nagaoka's model, having a compact charge, would agree with the scattering data. The Saturnian model had previously been rejected on other grounds. The so-called Rutherford model of the atom with orbiting electrons was not proposed by Rutherford in the 1911 paper.
Confirming the scattering theory: the 1913 experiment
In a 1913 paper, The Laws of Deflexion of α Particles through Large Angles, Geiger and Marsden describe a series of experiments by which they sought to experimentally verify Rutherford's equation. Rutherford's equation predicted that the number of scintillations per minute s that will be observed at a given angle Φ should be proportional to:
cosec4
thickness of foil t
magnitude of the square of central charge Qn
Their 1913 paper describes four experiments by which they proved each of these four relationships.
To test how the scattering varied with the angle of deflection (i.e. if s ∝ csc4). Geiger and Marsden built an apparatus that consisted of a hollow metal cylinder mounted on a turntable. Inside the cylinder was a metal foil (F) and a radiation source containing radon (R), mounted on a detached column (T) which allowed the cylinder to rotate independently. The column was also a tube by which air was pumped out of the cylinder. A microscope (M) with its objective lens covered by a fluorescent zinc sulfide screen (S) penetrated the wall of the cylinder and pointed at the metal foil. They tested with silver and gold foils. By turning the table, the microscope could be moved a full circle around the foil, allowing Geiger to observe and count alpha particles deflected by up to 150°. Correcting for experimental error, Geiger and Marsden found that the number of alpha particles that are deflected by a given angle Φ is indeed proportional to csc4.
Geiger and Marsden then tested how the scattering varied with the thickness of the foil (i.e. if s ∝ t). They constructed a disc (S) with six holes drilled in it. The holes were covered with metal foil (F) of varying thickness, or none for control. This disc was then sealed in a brass ring (A) between two glass plates (B and C). The disc could be rotated by means of a rod (P) to bring each window in front of the alpha particle source (R). On the rear glass pane was a zinc sulfide screen (Z). Geiger and Marsden found that the number of scintillations that appeared on the screen was indeed proportional to the thickness, as long as the thickness was small.
Geiger and Marsden reused the apparatus to measure how the scattering pattern varied with the square of the nuclear charge (i.e. if s ∝ Qn2). Geiger and Marsden did not know what the positive charge of the nucleus of their metals were (they had only just discovered the nucleus existed at all), but they assumed it was proportional to the atomic weight, so they tested whether the scattering was proportional to the atomic weight squared. Geiger and Marsden covered the holes of the disc with foils of gold, tin, silver, copper, and aluminium. They measured each foil's stopping power by equating it to an equivalent thickness of air. They counted the number of scintillations per minute that each foil produced on the screen. They divided the number of scintillations per minute by the respective foil's air equivalent, then divided again by the square root of the atomic weight (Geiger and Marsden knew that for foils of equal stopping power, the number of atoms per unit area is proportional to the square root of the atomic weight). Thus, for each metal, Geiger and Marsden obtained the number of scintillations that a fixed number of atoms produce. For each metal, they then divided this number by the square of the atomic weight, and found that the ratios were about the same. Thus they proved that s ∝ Qn2.
Finally, Geiger and Marsden tested how the scattering varied with the velocity of the alpha particles (i.e. if s ∝ ). Using the same apparatus, they slowed the alpha particles by placing extra sheets of mica in front of the alpha particle source. They found that, within the range of experimental error, the number of scintillations was indeed proportional to .
Positive charge on nucleus: 1913
In his 1911 paper (see above), Rutherford assumed that the central charge of the atom was positive, but a negative charge would have fitted his scattering model just as well. In a 1913 paper, Rutherford declared that the "nucleus" (as he now called it) was indeed positively charged, based on the result of experiments exploring the scattering of alpha particles in various gases.
In 1917, Rutherford and his assistant William Kay began exploring the passage of alpha particles through gases such as hydrogen and nitrogen. In this experiment, they shot a beam of alpha particles through hydrogen, and they carefully placed their detector—a zinc sulfide screen—just beyond the range of the alpha particles, which were absorbed by the gas. They nonetheless picked up charged particles of some sort causing scintillations on the screen. Rutherford interpreted this as alpha particles knocking the hydrogen nuclei forwards in the direction of the beam, not backwards.
Rutherford's scattering model
Rutherford begins his 1911 paper with a discussion of Thomson's results on scattering of beta particles, a form of radioactivity that results in high velocity electrons. Thomson's model had electrons circulating inside of a sphere of positive charge. Rutherford highlights the need for compound or multiple scattering events: the deflections predicted for each collision are much less than one degree. He then proposes a model which will produce large deflections on a single encounter: place all of the positive charge at the centre of the sphere and ignore the electron scattering as insignificant. The concentrated charge will explain why most alpha particles do not scatter to any measurable degree – they fly past too far from the charge – and yet particles that do pass very close to the centre scatter through large angles.
Maximum nuclear size estimate
Rutherford begins his analysis by considering a head-on collision between the alpha particle and atom. This will establish the minimum distance between them, a value which will be used throughout his calculations.
Assuming there are no external forces and that initially the alpha particles are far from the nucleus, the inverse-square law between the charges on the alpha particle and nucleus gives the potential energy gained by the particle as it approaches the nucleus. For head-on collisions between alpha particles and the nucleus, all the kinetic energy of the alpha particle is turned into potential energy and the particle stops and turns back.
Where the particle stops at a distance from the centre, the potential energy matches the original kinetic energy:
where
Rearranging:
For an alpha particle:
(mass) = =
(for the alpha particle) = 2 × =
(for gold) = 79 × =
(initial velocity) = (for this example)
The distance from the alpha particle to the centre of the nucleus () at this point is an upper limit for the nuclear radius.
Substituting these in gives the value of about , or 27 fm. (The true radius is about 7.3 fm.) The true radius of the nucleus is not recovered in these experiments because the alphas do not have enough energy to penetrate to more than 27 fm of the nuclear centre, as noted, when the actual radius of gold is 7.3 fm.
Rutherford's 1911 paper started with a slightly different formula suitable for head-on collision with a sphere of positive charge:
In Rutherford's notation, e is the elementary charge, N is the charge number of the nucleus (now also known as the atomic number), and E is the charge of an alpha particle. The convention in Rutherford's time was to measure charge in electrostatic units, distance in centimeters, force in dynes, and energy in ergs. The modern convention is to measure charge in coulombs, distance in meters, force in newtons, and energy in joules. Using coulombs requires using the Coulomb constant (k) in the equation. Rutherford used b as the turning point distance (called rmin above) and R is the radius of the atom. The first term is the Coulomb repulsion used above. This form assumes the alpha particle could penetrate the positive charge. At the time of Rutherford's paper, Thomson's plum pudding model proposed a positive charge with the radius of an atom, thousands of times larger than the rmin found above. Figure 1 shows how concentrated this potential is compared to the size of the atom.
Many of Rutherford's results are expressed in terms of this turning point distance rmin, simplifying the results and limiting the need for units to this calculation of turning point.
Single scattering by a heavy nucleus
From his results for a head on collision, Rutherford knows that alpha particle scattering occurs close to the centre of an atom, at a radius 10,000 times smaller than the atom. The electrons have negligible effect. He begins by assuming no energy loss in the collision, that is he ignores the recoil of the target atom. He will revisit each of these issues later in his paper.
Under these conditions, the alpha particle and atom interact through a central force, a physical problem studied first by Isaac Newton.
A central force only acts along a line between the particles and when the force varies with the inverse square, like Coulomb force in this case, a detailed theory was developed under the name of the Kepler problem. The well-known solutions to the Kepler problem are called orbits and unbound orbits are hyperbolas.
Thus Rutherford proposed that the alpha particle will take a hyperbolic trajectory in the repulsive force near the centre of the atom as shown in Figure 2.
To apply the hyperbolic trajectory solutions to the alpha particle problem, Rutherford expresses the parameters of the hyperbola in terms of the scattering geometry and energies. He starts with conservation of angular momentum. When the particle of mass and initial velocity is far from the atom, its angular momentum around the centre of the atom will be where is the impact parameter, which is the lateral distance between the alpha particle's path and the atom. At the point of closest approach, labeled A in Figure 2, the angular momentum will be . Therefore
Rutherford also applies the law of conservation of energy between the same two points:
The left hand side and the first term on the right hand side are the kinetic energies of the particle at the two points; the last term is the potential energy due to the Coulomb force between the alpha particle and atom at the point of closest approach (A). qa is the charge of the alpha particle, qg is the charge of the nucleus, and k is the Coulomb constant.
The energy equation can then be rearranged thus:
For convenience, the non-geometric physical variables in this equation can be contained in a variable , which is the point of closest approach in a head-on collision scenario which was explored in a previous section of this article:
This allows Rutherford simplify the energy equation to:
This leaves two simultaneous equations for , the first derived from the conservation of momentum equation and the second from the conservation of energy equation. Eliminating and gives at a new formula for :
The next step is to find a formula for . From Figure 2, is the sum of two distances related to the hyperbola, SO and OA. Using the following logic, these distances can be expressed in terms of angle and impact parameter .
The eccentricity of a hyperbola is a value that describes the hyperbola's shape. It can be calculated by dividing the focal distance by the length of the semi-major axis, which per Figure 2 is . As can be seen in Figure 3, the eccentricity is also equal to , where is the angle between the major axis and the asymptote. Therefore:
As can be deduced from Figure 2, the focal distance SO is
and therefore
With these formulas for SO and OA, the distance can be written in terms of and simplified using a trigonometric identity known as a half-angle formula:
Applying a trigonometric identity known as the cotangent double angle formula and the previous equation for gives a simpler relationship between the physical and geometric variables:
The scattering angle of the particle is and therefore . With the help of a trigonometric identity known as a reflection formula, the relationship between θ and b can be resolved to:
which can be rearranged to give
Rutherford gives some illustrative values as shown in this table:
Rutherford's approach to this scattering problem remains a standard treatment in textbooks on classical mechanics.
Intensity vs angle
To compare to experiments the relationship between impact parameter and scattering angle needs to be converted to probability versus angle. The scattering cross section gives the relative intensity by angles:
In classical mechanics, the scattering angle is uniquely determined the initial kinetic energy of the incoming particles and the impact parameter . Therefore, the number of particles scattered into an angle between and must be the same as the number of particles with associated impact parameters between and . For an incident intensity , this implies:
Thus the cross section depends on scattering angle as:
Using the impact parameter as a function of angle, , from the single scattering result above produces the Rutherford scattering cross section:
s = the number of alpha particles falling on unit area at an angle of deflection Φ
r = distance from point of incidence of α rays on scattering material
X = total number of particles falling on the scattering material
n = number of atoms in a unit volume of the material
t = thickness of the foil
qn = positive charge of the atomic nucleus
qa = positive charge of the alpha particles
m = mass of an alpha particle
v = velocity of the alpha particle
This formula predicted the results that Geiger measured in the coming year. The scattering probability into small angles greatly exceeds the probability in to larger angles, reflecting the tiny nucleus surrounded by empty space. However, for rare close encounters, large angle scattering occurs with just a single target.
At the end of his development of the cross section formula, Rutherford emphasises that the results apply to single scattering and thus require measurements with thin foils. For thin foils the degree of scattering is proportional to the foil thickness in agreement with Geiger's measurements.
Comparison to JJ Thomson's results
At the time of Rutherford's paper, JJ Thomson was the "undisputed world master in the design of atoms". Rutherford needed to compare his new approach to Thomson's. Thomson's model, presented in 1910, modelled the electron collisions with hyperbolic orbits from his 1906 paper combined with a factor for the positive sphere. Multiple resulting small deflections compounded using a random walk.
In his paper Rutherford emphasised that single scattering alone could account for Thomson's results if the positive charge were concentrated in the centre.
Rutherford computes the probability of single scattering from a compact charge and demonstrates that it is 3 times larger than Thomson's multiple scattering probability. Rutherford completes his analysis including the effects of density and foil thickness, then concludes that thin foils are governed by single scattering, not multiple scattering.
Later analysis showed Thomson's scattering model could not account for large scattering.
The maximum angular deflection from electron scattering or from the positive sphere each come to less than 0.02°; even many such scattering events compounded would result in less than a one degree average deflection and a probability of scattering through 90° of less than one in 103500.
Target recoil
Rutherford's analysis assumed that alpha particle trajectories turned at the centre of the atom but the exit velocity was not reduced. This is equivalent to assuming that the concentrated charge at the centre had infinite mass or was anchored in place. Rutherford discusses the limitations of this assumption by comparing scattering from lighter atoms like aluminium with heavier atoms like gold. If the concentrated charge is lighter it will recoil from the interaction, gaining momentum while the alpha particle loses momentum and consequently slows down.
Modern treatments analyze this type of Coulomb scattering in the centre of mass reference frame. The six coordinates of the two particles (also called "bodies") are converted into three relative coordinates between the two particles and three centre-of-mass coordinates moving in space (called the lab frame). The interaction only occurs in the relative coordinates, giving an equivalent one-body problem just as Rutherford solved, but with different interpretations for the mass and scattering angle.
Rather than the mass of the alpha particle, the more accurate formula including recoil uses reduced mass:
For Rutherford's alpha particle scattering from gold, with mass of 197, the reduced mass is very close to the mass of the alpha particle:
For lighter aluminium, with mass 27, the effect is greater:
a 13% difference in mass. Rutherford notes this difference and suggests experiments be performed with lighter atoms.
The second effect is a change in scattering angle. The angle in the relative coordinate system or centre of mass frame needs to be converted to an angle in the lab frame. In the lab frame, denoted by a subscript L, the scattering angle for a general central potential is
For a heavy particle like gold used by Rutherford, the factor can be neglected at almost all angles. Then the lab and relative angles are the same, .
The change in scattering angle alters the formula for differential cross-section needed for comparison to experiment. For any central potential, the differential cross-section in the lab frame is related to that in the centre-of-mass frame by
where
Limitations to Rutherford's scattering formula
Very light nuclei and higher energies
In 1919 Rutherford analyzed alpha particle scattering from hydrogen atoms, showing the limits of the 1911 formula even with corrections for reduced mass. Similar issues with smaller deviations for helium, magnesium, aluminium lead to the conclusion that the alpha particle was penetrating the nucleus in these cases. This allowed the first estimates of the size of atomic nuclei. Later experiments based on cyclotron acceleration of alpha particles striking heavier nuclei provided data for analysis of interaction between the alpha particle and the nuclear surface. However at energies that push the alpha particles deeper they are strongly absorbed by the nuclei, a more complex interaction.
Quantum mechanics
Rutherford's treatment of alpha particle scattering seems to rely on classical mechanics and yet the particles are of sub-atomic dimensions. However the critical aspects of the theory ultimately rely on conservation of momentum and energy. These concepts apply equally in classical and quantum regimes: the scattering ideas developed by Rutherford apply to subatomic elastic scattering problems like neutron-proton scattering.
An alternative method to find the scattering angle
This section presents an alternative method to find the relation between the impact parameter and deflection angle in a single-atom encounter, using a force-centric approach as opposed to the energy-centric one that Rutherford used.
The scattering geometry is shown in this diagram
The impact parameter b is the distance between the alpha particle's initial trajectory and a parallel line that goes through the nucleus. Smaller values of b bring the particle closer to the atom so it feels more deflection force resulting in a larger deflection angle θ. The goal is to find the relationship between b and the deflection angle.
The alpha particle's path is a hyperbola and the net change in momentum runs along the axis of symmetry. From the geometry in the diagram and the magnitude of the initial and final momentum vectors, , the magnitude of can be related to the deflection angle:
A second formula for involving b will give the relationship to the deflection angle.
The net change in momentum can also be found by adding small increments to momentum all along the trajectory using the integral
where is the distance between the alpha particle and the centre of the nucleus and is its angle from the axis of symmetry. These two are the polar coordinates of the alpha particle at time .
Here the Coulomb force exerted along the line between the alpha particle and the atom is and the factor gives that part of the force causing deflection.
The polar coordinates r and φ depend on t in the integral, but they must be related to each other as they both vary as the particle moves. Changing the variable and limits of integration from t to φ makes this connection explicit:
The factor is the reciprocal of the angular velocity the particle. Since the force is only along the line between the particle and the atom, the angular momentum, which is proportional to the angular velocity, is constant:
This law of conservation of angular momentum gives a formula for :
Replacing in the integral for ΔP simultaneously eliminates the dependence on r:
Applying the trigonometric identities and to simplify this result gives the second formula for :
Solving for θ as a function of b gives the final result
Why the plum pudding model was wrong
J. J. Thomson himself didn't study alpha particle scattering, but he did study beta particle scattering. In his 1910 paper "On the Scattering of rapidly moving Electrified Particles", Thomson presented equations that modelled how beta particles scatter in a collision with an atom. Rutherford adapted those equations to alpha particle scattering in his 1911 paper "The Scattering of α and β Particles by Matter and the Structure of the Atom".
Deflection by the positive sphere
In Thomson's 1910 paper "On the Scattering of rapidly moving Electrified Particles", Thomson presented the following equation (in this article's notation) that isolates the effect of the positive sphere in the plum pudding model on an incoming beta particle.
Thomson did not explain how he arrived at this equation, but this section provides an educated guess and at the same time adapts the equation to alpha particle scattering.
Consider an alpha particle passing by a positive sphere of pure positive charge (no electrons) with a radius R and mass equal to those of a gold atom. The alpha particle passes just close enough to graze the edge of the sphere, which is where the electric field of the sphere is strongest.
An earlier section of this article presented an equation which models how an incoming charged particle is deflected by another charged particle at a fixed position.
This equation can be used to calculate the deflection angle in the special case in Figure 4 by setting the impact parameter b to the same value as the radius of the sphere R. So long as the alpha particle does not penetrate the sphere, there is no practical difference between a sphere of charge and a point charge.
qg = positive charge of the gold atom = =
qa = charge of the alpha particle = =
R = radius of the gold atom =
v = speed of the alpha particle =
m = mass of the alpha particle =
k = Coulomb constant =
This shows that the largest possible deflection will be very small, to the point that the path of the alpha particle passing through the positive sphere of a gold atom is almost a straight line. Therefore in computing the average deflection, which will be smaller still, we will treat the particle's path through the sphere as a chord of length L.
Inside a sphere of uniformly distributed positive charge, the force exerted on the alpha particle at any point along its path through the sphere is
The lateral component of this force is
The lateral change in momentum py is therefore
The deflection angle is given by
where px is the average horizontal momentum, which is first reduced then restored as horizontal force changes direction as the alpha particle goes across the sphere. Since the deflection is very small, can be treated as equal to . The chord length , per Pythagorean theorem.
The average deflection angle sums the angle for values of b and L across the entire sphere and divides by the cross-section of the sphere:
This matches Thomson's formula in his 1910 paper.
Deflection by the electrons
Consider an alpha particle passing through an atom of radius R along a path of length L. The effect of the positive sphere is ignored so as to isolate the effect of the atomic electrons. As with the positive sphere, deflection by the electrons is expected to be very small, to the point that the path is practically a straight line.
For the electrons within an arbitrary distance s of the alpha particle's path, their mean distance will be s. Therefore, the average deflection per electron will be
where qe is the elementary charge. The average net deflection by all the electrons within this arbitrary cylinder of effect around the alpha particle's path is
where N0 is the number of electrons per unit volume and is the volume of this cylinder.
Treating L as a straight line, where b is the distance of this line from the centre. The mean of is therefore
To obtain the mean deflection , replace in the equation for :
where N is the number of electrons in the atom, equal to .
Cumulative effect
Applying Thomson's equations described above to an alpha particle colliding with a gold atom, using the following values:
qg = positive charge of the gold atom = =
qa = charge of the alpha particle = =
qe = elementary charge =
R = radius of the gold atom =
v = speed of the alpha particle =
m = mass of the alpha particle =
k = Coulomb constant =
N = number of electrons in the gold atom = 79
gives the average angle by which the alpha particle should be deflected by the atomic electrons as:
The average angle by which an alpha particle should be deflected by the positive sphere is:
The net deflection for a single atomic collision is:
On average the positive sphere and the electrons alike provide very little deflection in a single collision. Thomson's model combined many single-scattering events from the atom's electrons and a positive sphere. Each collision may increase or decrease the total scattering angle. Only very rarely would a series of collisions all line up in the same direction. The result is similar to the standard statistical problem called a random walk. If the average deflection angle of the alpha particle in a single collision with an atom is , then the average deflection after n collisions is
The probability that an alpha particle will be deflected by a total of more than 90° after n deflections is given by:
where e is Euler's number (≈2.71828...). A gold foil with a thickness of 1.5 micrometers would be about 10,000 atoms thick. If the average deflection per atom is 0.008°, the average deflection after 10,000 collisions would be 0.8°. The probability of an alpha particle being deflected by more than 90° will be
While in Thomson's plum pudding model it is mathematically possible that an alpha particle could be deflected by more than 90° after 10,000 collisions, the probability of such an event is so low as to be undetectable. This extremely small number shows that Thomson's model cannot explain the results of the Geiger-Mardsen experiment of 1909.
| Physical sciences | Atomic physics | null |
320217 | https://en.wikipedia.org/wiki/Substation | Substation | A substation is a part of an electrical generation, transmission, and distribution system. Substations transform voltage from high to low, or the reverse, or perform any of several other important functions. Between the generating station and consumer, electric power may flow through several substations at different voltage levels. A substation may include transformers to change voltage levels between high transmission voltages and lower distribution voltages, or at the interconnection of two different transmission voltages. They are a common component of the infrastructure. There are 55,000 substations in the United States.
Substations may be owned and operated by an electrical utility, or may be owned by a large industrial or commercial customer. Generally substations are unattended, relying on SCADA for remote supervision and control.
The word substation comes from the days before the distribution system became a grid. As central generation stations became larger, smaller generating plants were converted to distribution stations, receiving their energy supply from a larger plant instead of using their own generators. The first substations were connected to only one power station, where the generators were housed, and were subsidiaries of that power station.
Construction
Substations may be designed and built by a contractor or alternately all phases of its development may be handled by the electrical utility. Most commonly, the utility does the engineering and procurement while hiring a contractor for actual construction. Major design constraints for construction of substations include land availability and cost, limitations on the construction period, transportation restrictions, and the need to get the substation running quickly. Prefabrication is a common way to reduce the construction cost. For connecting the new substation, a partial outage at another substation may be required, but the utility often tries to minimize downtime.
Types
Substations typically serve at least one of the following purposes:
Increasing the voltage produced by electric power generation for efficient transmission over long distances, using step-up transformers
Interconnection of different power grids
Reducing the voltage from transmission to lower-voltage distribution lines that supply individual homes or businesses
Converting from alternating current (AC) to direct current (DC)
Transmission substation
A transmission substation connects two or more transmission lines. The simplest case is where all transmission lines have the same voltage. In such cases, substation contains high-voltage switches that allow lines to be connected or isolated for fault clearance or maintenance. A transmission station may have transformers to convert between two transmission voltages, voltage control/power factor correction devices such as capacitors, reactors or static VAR compensators and equipment such as phase shifting transformers to control power flow between two adjacent power systems.
Transmission substations can range from simple to complex. A small "switching station" may be little more than a bus plus some circuit breakers. The largest transmission substations can cover a large area (several acres/hectares) with multiple voltage levels, many circuit breakers, and a large amount of protection and control equipment (voltage and current transformers, relays and SCADA systems). Modern substations may be implemented using international standards such as IEC Standard 61850.
Distribution substation
A distribution substation transfers power from the transmission system to the distribution system of an area. It is uneconomical to directly connect electricity consumers to the main transmission network, unless they use large amounts of power, so the distribution station reduces voltage to a level suitable for local distribution.
The input for a distribution substation is typically at least two transmission or sub-transmission lines. Input voltage may be, for example, 115 kV, or whatever is common in the area. The output is a number of feeders. Distribution voltages are typically medium voltage, between 2.4 kV and 33 kV, depending on the size of the area served and the practices of the local utility. The feeders run along streets overhead (or underground, in some cases) and power the distribution transformers at or near the customer premises.
In addition to transforming voltage, distribution substations also isolate faults in either the transmission or distribution systems. Distribution substations are typically the points of voltage regulation, although on long distribution circuits (of several miles/kilometers), voltage regulation equipment may also be installed along the line.
The downtown areas of large cities feature complicated distribution substations, with high-voltage switching, and switching and backup systems on the low-voltage side. More typical distribution substations have a switch, one transformer, and minimal facilities on the low-voltage side.
Collector substation
In distributed generation projects such as a wind farm or photovoltaic power station, a collector substation may be required. It resembles a distribution substation although power flow is in the opposite direction, from many wind turbines or inverters up into the transmission grid. Usually for economy of construction the collector system operates around 35 kV, although some collector systems are 12 kV, and the collector substation steps up voltage to a transmission voltage for the grid. The collector substation can also provide power factor correction if it is needed, metering, and control of the wind farm. In some special cases a collector substation can also contain an HVDC converter station.
Collector substations also exist where multiple thermal or hydroelectric power plants of comparable output power are in proximity. Examples for such substations are Brauweiler in Germany and Hradec in the Czech Republic, where power is collected from nearby lignite-fired power plants. If no transformers are required for increasing the voltage to transmission level, the substation is a switching station.
Converter substations
Converter substations may be associated with HVDC converter plants, traction current, or interconnected non-synchronous networks. These stations contain power electronic devices to change the frequency of current, or else convert from alternating to direct current or the reverse. Formerly rotary converters changed frequency to interconnect two systems; nowadays such substations are rare.
Switching station
A switching station is a substation without transformers and operating only at a single voltage level. Switching stations are sometimes used as collector and distribution stations. Sometimes they are used for switching the current to back-up lines or for parallelizing circuits in case of failure. An example is the switching stations for the HVDC Inga–Shaba transmission line.
A switching station may also be known as a switchyard, and these are commonly located directly adjacent to or nearby a power station. In this case the generators from the power station supply their power into the yard onto the generator bus on one side of the yard, and the transmission lines take their power from a Feeder Bus on the other side of the yard.
An important function performed by a substation is switching, which is the connecting and disconnecting of transmission lines or other components to and from the system. Switching events may be planned or unplanned. A transmission line or other component may need to be de-energized for maintenance or for new construction, for example, adding or removing a transmission line or a transformer. To maintain reliability of supply, companies aim at keeping the system up and running while performing maintenance. All work to be performed, from routine testing to adding entirely new substations, should be done while keeping the whole system running.
Unplanned switching events are caused by a fault in a transmission line or any other component, for example:
a line is hit by lightning and develops an arc,
a tower is blown down by high wind.
The function of the switching station is to isolate the faulty portion of the system in the shortest possible time. De-energizing faulty equipment protects it from further damage, and isolating a fault helps keep the rest of the electrical grid operating with stability.
Railways
Electrified railways also use substations, often distribution substations. In some cases a conversion of the current type takes place, commonly with rectifiers for direct current (DC) trains, or rotary converters for trains using alternating current (AC) at frequencies other than that of the public grid. Sometimes they are also transmission substations or collector substations if the railway network also operates its own grid and generators to supply the other stations.
Mobile substation
A mobile substation is a substation on wheels, containing a transformer, breakers and buswork mounted on a self-contained semi-trailer, meant to be pulled by a truck. They are designed to be compact for travel on public roads, and are used for temporary backup in times of natural disaster or war. Mobile substations are usually rated much lower than permanent installations, and may be built in several units to meet road travel limitations.
Design
Substation design is aimed at minimizing cost while ensuring power availability and reliability, and enabling changes to the substation in the future.
Substations may be built outdoors, indoors, or underground or in a combination of these locations.
Location selection
Selection of the location of a substation must consider many factors. Sufficient land area is required for installation of equipment with necessary clearances for electrical safety, and for access to maintain large apparatus such as transformers. The site must have room for expansion due to load growth or planned transmission additions. Environmental effects of the substation must be considered, such as drainage, noise and road traffic effects.
The substation site must be reasonably central to the distribution area to be served. The site must be secure from intrusion by passers-by, both to protect people from injury by electric shock or arcs, and to protect the electrical system from misoperation due to vandalism.
If not owned and operated by a utility company, substations are typically occupied on a long lease such as a renewable 99-year lease, giving the utility company security of tenure.
Design diagrams
The first step in planning a substation layout is the preparation of a one-line diagram, which shows in simplified form the switching and protection arrangement required, as well as the incoming supply lines and outgoing feeders or transmission lines. It is a usual practice by many electrical utilities to prepare one-line diagrams with principal elements (lines, switches, circuit breakers, transformers) arranged on the page similarly to the way the apparatus would be laid out in the actual station.
In a common design, incoming lines have a disconnector and a circuit breaker. In some cases, the lines will not have both, with either a switch or a circuit breaker being all that is considered necessary. A disconnect switch is used to provide isolation, since it cannot interrupt load current. A circuit breaker is used as a protection device to interrupt fault currents automatically, and may be used to switch loads on and off, or to cut off a line when power is flowing in the 'wrong' direction. When a large fault current flows through the circuit breaker, this is detected through the use of current transformers. The magnitude of the current transformer outputs may be used to trip the circuit breaker resulting in a disconnection of the load supplied by the circuit break from the feeding point. This seeks to isolate the fault point from the rest of the system, and allow the rest of the system to continue operating with minimal impact. Both switches and circuit breakers may be operated locally (within the substation) or remotely from a supervisory control center.
With overhead transmission lines, the propagation of lightning and switching surges can cause insulation failures into substation equipment. Line entrance surge arrestors are used to protect substation equipment accordingly. Insulation Coordination studies are carried out extensively to ensure equipment failure (and associated outages) is minimal.
Once past the switching components, the lines of a given voltage connect to one or more buses. These are sets of busbars, usually in multiples of three, since three-phase electrical power distribution is largely universal around the world.
The arrangement of switches, circuit breakers, and buses used affects the cost and reliability of the substation. For important substations a ring bus, double bus, or so-called "breaker and a half" setup can be used, so that the failure of any one circuit breaker does not interrupt power to other circuits, and so that parts of the substation may be de-energized for maintenance and repairs. Substations feeding only a single industrial load may have minimal switching provisions, especially for small installations.
Safety
Because of the risk of electrical shock, substations are inherently dangerous to electrical workers. To mitigate this hazard, substations are designed with various safety features. Live conductors and bare equipment are kept separate, either with protected equipment, or using screens or distance. Based on the jurisdiction or company, there are safety standards with minimum required clearance between different live equipment or conductors or between live metal and the ground, which often varies with higher clearance being required for higher voltages because of the greater ability to generate flashover. To this is added the necessary space for employees to work safely and vehicles to pass. Sometimes it is necessary to work on parts of the substation while energized, but employees must maintain a safe distance of at least . The aim to reduce substation footprints comes into conflict with ease of maintenance enhanced by including gaps where employees can safely work.
Underneath a substation, a mat or grid of conductors laid around underground provides grounding. This grid, which is typically copper although it may be galvanized iron in some countries, is used to ground circuits that are being worked on to prevent accidental re-energization while workers are in contact with a de-energized circuit. Often, earth rods are driven deeper into the ground from the grounding grid for lower resistance grounding, and may be surrounded by bentonite or marconite to further reduce resistance and ensure effective grounding for the lifetime of the substation. Above ground, the grounding conductors may be steel, aluminum, or copper. They must be thick enough to carry the expected current of a fault for 1-3 seconds and remain undamaged. Substation fences, typically at least in height, both protect the public from electrical hazards and also protect the substation from vandalism. Internal fences can also be incorporated to protect employees from areas that are unsafe when energized.
Components
Substations generally have switching, protection and control equipment, and transformers. In a large substation, circuit breakers are used to interrupt any short circuits or overload currents that may occur on the network. Smaller distribution stations may use recloser circuit breakers or fuses for protection of distribution circuits. Substations themselves do not usually have generators, although a power plant may have a substation nearby. Other devices such as capacitors, voltage regulators, and reactors may also be located at a substation.
Substations may be on the surface in fenced enclosures, underground, or special-purpose buildings. High-rise buildings may have several indoor substations. Indoor substations are usually found in urban areas to reduce the noise from transformers, improve appearance, or protect switchgear from extreme climate or pollution.
Substations often use busbars as conductors between electrical equipment. Busbars may be aluminum tubing thick, or else wires (strain bus).
Outdoor, above-ground substation structures include wood pole, lattice metal tower, and tubular metal structures, although other variants are available. Where space is plentiful and appearance of the station is not a factor, steel lattice towers provide low-cost supports for transmission lines and apparatus. Low-profile substations may be specified in suburban areas where appearance is more critical. Indoor substations may be gas insulated substations (GIS) (at high voltages, with gas insulated switchgear), or use metal-enclosed or metal-clad switchgear at lower voltages. Urban and suburban indoor substations may be finished on the outside so as to blend in with other buildings in the area.
A compact substation is generally an outdoor substation built in a metal enclosure, in which each item of the electrical equipment is located very near to each other to create a relatively smaller footprint size of the substation.
Switchgear
High-voltage circuit breakers are commonly used to interrupt the flow of current in substation equipment. At the time of interruption, current could be normal, too high due to excessive load, unusual due to a fault, or tripped by protective relays prior to anticipated trouble. The most common technologies to extinguish the power arc from separating the conductors in the breaker include:
Air at atmospheric pressure (air-insulated switchgear (AIS)), which is the most common worldwide. Air is the cheapest insulator and is easy to modify, but AIS takes up more space, and leaves equipment exposed to the outside environment. One drawback of AIS is the visual impact of a larger substation with overhead power lines entering and exiting, which may be unacceptable in scenic or urban areas. AIS requires additional bracing in a seismically active area, and emits more electromagnetic fields and noise than alternative technologies.
Gas (gas circuit breaker (GCB) or gas-insulated switchgear (GIS)), most commonly sulfur hexafluoride (SF6) or a mixture of gases including SF6. Although it is the most expensive, these gases are a much more effective insulator than air. GIS require only 10 to 20 percent of the land area as AIS, which can save on land acquisition cost in urban areas, and allow the substation to be built at the exact location where its power is being used in an industrial or urban area—which can be a significant cost savings. On the generation side, GIS can be installed closer to the generator which allows cost savings in cabling, bus duct connections, and civil construction and can increase reliability. GIS can replace AIS if power requirements increase without requiring additional land area. Additionally, GIS is commonly installed in an enclosed building that keeps the equipment protected from pollution and salt. Unless the substation is often used for switching, maintenance cost can be very low or even zero for many years. Because SF6 turns to solid around , in some climates these circuit breakers require heaters to function in extremely cold weather. SF6 has been used in switchgear since the 1960s.
Mineral oil (called OCB for oil circuit breaker) provides a high resistance between the opened contacts, effectively stopping the flow of current. Although oil circuit breakers are suitable for a wide range of voltages, the oil becomes contaminated during the suppression of arcs and must be filtered or replaced periodically.
Vacuum is a better insulator than air but less than gas or oil. Vacuum circuit breakers (VCB) are smaller than air circuit breakers and are commonly used in distribution and other switchgear under 35kv.
Mixed, including both gas and air insulation. Although it’s the least common option it can be useful when an air-insulated substation needs to be expanded but there is very limited location for additional construction.
Reclosers are similar to breakers, and can be cheaper because they do not require separate protective relays. Often used in distribution, they often are programmed to trip when the amps exceed a certain amount over a period of time. Reclosers will attempt to re-energize the circuit after a delay. If unsuccessful for a few times, the recloser will have to be manually reset by an electrical worker.
Capacitors
Capacitor banks are used in substations to balance the lagging current draw from inductive loads (such as motors, transformers, and some industrial equipment) with their reactive load. Additional capacitor capacity may be needed if dispersed generation (such as small diesel generators, rooftop photovoltaic solar panels, or wind turbines) are added to the system. Capacitors can reduce the current in wires, helping stem system losses from voltage drop or enabling extra power to be sent through the conductors. Capacitors may be left on in response to constant inductive load or turned on when inductive load is increased, such as in the summer for air conditioners. The switching may be remote and can be done manually or automatically.
Control rooms
Larger substations have control rooms for the equipment used to monitor, control, and protect the rest of the substation equipment. It often contains protective relays, meters, breaker controls, communications, batteries, and recorders that save detailed data about substation operations, particularly when there is any unusual activity, to help reconstruct what happened after the fact. These control rooms typically are heated and air conditioned to ensure the reliable operation of this equipment. Additional equipment is necessary to handle power surges associated with intermittent renewable energy such as dispersed generation from wind or solar.
Transformers
Most transformers lose between 5 and 1.5 percent of their input as heat and noise. Iron losses are no-load and constant whenever the transformer is energized, while copper and auxiliary losses are proportionate to the square of the current. Auxiliary losses are due to running fans and pumps which is noisy when the transformer is operating at maximum capacity. To reduce noise, enclosures are often built around the transformer and can also be added after the substation is built.
Oil-based transformers are often built with bunded areas to prevent the escape of flaming or leaking oil. Fire separation areas or firewalls are built around the transformer to stop the spread of fire. Firefighting vehicles are allowed a path to access the area.
Maintenance
Maintenance of substations involves inspections, data collection and analysis, and routine scheduled work. Using methods such as infrared scanning and dissolved gas analysis, it can be predicted when the substation will need maintenance and predict dangers before they materialize. Infrared technology finds hot spots in the substation where electrical energy is being converted to heat, which indicates a problem and can cause additional damage from the high heat. Dissolved gas analysis can tell when an oil-insulated transformer needs to have the oil filtered or replace, and also detect other issues.
Automation
Early electrical substations required manual switching or adjustment of equipment, and manual collection of data for load, energy consumption, and abnormal events. As the complexity of distribution networks grew, it became economically necessary to automate supervision and control of substations from a centrally attended point, to allow overall coordination in case of emergencies and to reduce operating costs. Early efforts to remote control substations used dedicated communication wires, often run alongside power circuits. Power-line carrier, microwave radio, fiber optic cables as well as dedicated wired remote control circuits have all been applied to Supervisory Control and Data Acquisition (SCADA) for substations. The development of the microprocessor made for an exponential increase in the number of points that could be economically controlled and monitored. Today, standardized communication protocols such as DNP3, IEC 61850 and Modbus, to list a few, are used to allow multiple intelligent electronic devices to communicate with each other and supervisory control centers. Distributed automatic control at substations is one element of the so-called smart grid.
| Technology | Electricity transmission and distribution | null |
1506024 | https://en.wikipedia.org/wiki/Crystal%20violet | Crystal violet | Crystal violet or gentian violet, also known as methyl violet 10B or hexamethyl pararosaniline chloride, is a triarylmethane dye used as a histological stain and in Gram's method of classifying bacteria. Crystal violet has antibacterial, antifungal, and anthelmintic (vermicide) properties and was formerly important as a topical antiseptic. The medical use of the dye has been largely superseded by more modern drugs, although it is still listed by the World Health Organization.
The name gentian violet was originally used for a mixture of methyl pararosaniline dyes (methyl violet), but is now often considered a synonym for crystal violet. The name refers to its colour, being like that of the petals of certain gentian flowers; it is not made from gentians or violets.
Production
A number of possible routes can be used to prepare crystal violet. The original procedure developed by the German chemists Kern and Caro involved the reaction of dimethylaniline with phosgene to give 4,4′-bis(dimethylamino)benzophenone (Michler's ketone) as an intermediate. This was then reacted with additional dimethylaniline in the presence of phosphorus oxychloride and hydrochloric acid.
The dye can also be prepared by the condensation of formaldehyde and dimethylaniline to give a leuco dye:
CH2O + 3 C6H5N(CH3)2 → CH(C6H4N(CH3)2)3 + H2O
Second, this colourless compound is oxidized to the coloured cationic form (hereafter with oxygen, but a typical oxidizing agent is manganese dioxide, MnO2):
CH(C6H4N(CH3)2)3 + HCl + O2 → [C(C6H4N(CH3)2)3]Cl + H2O
Dye colour
When dissolved in water, the dye has a blue-violet colour with an absorbance maximum at 590 nm and an extinction coefficient of 87,000 M−1 cm−1. The colour of the dye depends on the acidity of the solution. At a pH of +1.0, the dye is green with absorption maxima at 420 nm and 620 nm, while in a strongly acidic solution (pH −1.0), the dye is yellow with an absorption maximum at 420 nm.
The different colours are a result of the different charged states of the dye molecule. In the yellow form, all three nitrogen atoms carry a positive charge, of which two are protonated, while the green colour corresponds to a form of the dye with two of the nitrogen atoms positively charged. At neutral pH, both extra protons are lost to the solution, leaving only one of the nitrogen atoms positive charged. The pKa for the loss of the two protons are approximately 1.15 and 1.8.
In alkaline solutions, nucleophilic hydroxyl ions attack the electrophilic central carbon to produce the colourless triphenylmethanol or carbinol form of the dye. Some triphenylmethanol is also formed under very acidic conditions when the positive charges on the nitrogen atoms lead to an enhancement of the electrophilic character of the central carbon, which allows the nucleophilic attack by water molecules. This effect produces a slight fading of the yellow colour.
Applications
Industry
Crystal violet is used as a textile and paper dye, and is a component of navy blue and black inks for printing, ball-point pens, and inkjet printers. Historically, it was the most common dye used in early duplication machines, such as the mimeograph and the ditto machine. It is sometimes used to colourize diverse products such as fertilizer, antifreeze, detergent, and leather. Marking blue, used to mark out pieces in metalworking, is composed of methylated spirits, shellac, and gentian violet.
Science
When conducting DNA gel electrophoresis, crystal violet can be used as a nontoxic DNA stain as an alternative to fluorescent, intercalating dyes such as ethidium bromide. Used in this manner, it may be either incorporated into the agarose gel or applied after the electrophoresis process is finished. Used at a 10 ppm concentration and allowed to stain a gel after electrophoresis for 30 minutes, it can detect as little as 16 ng of DNA. Through use of a methyl orange counterstain and a more complex staining method, sensitivity can be improved further to 8 ng of DNA. When crystal violet is used as an alternative to fluorescent stains, it is not necessary to use ultraviolet illumination; this has made crystal violet popular as a means of avoiding UV-induced DNA destruction when performing DNA cloning in vitro.
In biomedical research, crystal violet can be used to stain the nuclei of adherent cells. In this application, crystal violet works as an intercalating dye and allows the quantification of DNA which is proportional to the number of cells.
The dye is used as a histological stain, particularly in Gram staining for classifying bacteria.
In forensics, crystal violet was used to develop fingerprints.
Crystal violet is also used as a tissue stain in the preparation of light microscopy sections. In laboratory, solutions containing crystal violet and formalin are often used to simultaneously fix and stain cells grown in tissue culture to preserve them and make them easily visible, since most cells are colourless. It is also sometimes used as a cheap way to put identification markings on laboratory mice; since many strains of lab mice are albino, the purple colour stays on their fur for several weeks.
Crystal violet can be used as an alternative to Coomassie brilliant blue (CBB) in staining of proteins separated by SDS-PAGE, reportedly showing a 5x improved sensitivity vs CBB.
Medical
Gentian violet has antibacterial, antifungal, antihelminthic, antitrypanosomal, antiangiogenic, and antitumor properties. It is used medically for these properties, in particular for dentistry, and is also known as "pyoctanin" (or "pyoctanine"). It is commonly used for:
Marking the skin for surgery preparation and allergy testing;
Treating Candida albicans and related fungal infections, such as thrush, yeast infections, various types of tinea (ringworm, athlete's foot, jock itch);
Treating impetigo; it was used primarily before the advent of antibiotics, but still useful to persons who may be allergic to penicillin.
In resource-limited settings, gentian violet is used to manage burn wounds, inflammation of the umbilical cord stump (omphalitis) in the neonatal period, oral candidiasis in HIV-infected patients and mouth ulcers in children with measles.
In body piercing, gentian violet is commonly used to mark the location for placing piercings, including surface piercings.
Veterinary
Because of its antimicrobial activity, it is used to treat ich in fish. However, it usually is illegal to use in fish intended for human consumption.
History
Synthesis
Crystal violet is one of the components of methyl violet, a dye first synthesized by Charles Lauth in 1861. From 1866, methyl violet was manufactured by the Saint-Denis-based firm of Poirrier et Chappat and marketed under the name "Violet de Paris". It was a mixture of the tetra-, penta- and hexamethylated pararosanilines.
Crystal violet itself was first synthesized in 1883 by (1850–1893) working in Basel at the firm of Bindschedler & Busch. To optimize the difficult synthesis which used the highly toxic phosgene, Kern entered into a collaboration with the German chemist Heinrich Caro at BASF. Kern also found that by starting with diethylaniline rather than dimethylaniline, he could synthesize the closely related violet dye now known as C.I. 42600 or C.I. Basic violet 4.
Gentian violet
The name "gentian violet" (or Gentianaviolett in German) is thought to have been introduced by the German pharmacist Georg Grübler, who in 1880 started a company in Leipzig that specialized in the sale of staining reagents for histology. The gentian violet stain marketed by Grübler probably contained a mixture of methylated pararosaniline dyes. The stain proved popular and in 1884 was used by Hans Christian Gram to stain bacteria. He credited Paul Ehrlich for the aniline-gentian violet mixture. Grübler's gentian violet was probably very similar, if not identical, to Lauth's methyl violet, which had been used as a stain by Victor André Cornil in 1875.
Although the name gentian violet continued to be used for the histological stain, the name was not used in the dye and textile industries. The composition of the stain was not defined and different suppliers used different mixtures. In 1922, the Biological Stain Commission appointed a committee chaired by Harold Conn to look into the suitability of the different commercial products. In his book Biological Stains, Conn describes gentian violet as a "poorly defined mixture of violet rosanilins".
The German ophthalmologist Jakob Stilling is credited with discovering the antiseptic properties of gentian violet. He published a monograph in 1890 on the bactericidal effects of a solution that he christened "pyoctanin", which was probably a mixture of aniline dyes similar to gentian violet. He set up a collaboration with E. Merck & Co. to market "Pyoktanin caeruleum" as an antiseptic.
In 1902, Drigalski and Conradi found that although crystal violet inhibited the growth of many bacteria, it has little effect on Bacillus coli (Escherichia coli) and Bacillus typhi (Salmonella typhi), which are both gram-negative bacteria. A much more detailed study of the effects of Grübler's gentian violet on different strains of bacteria was published by John Churchman in 1912. He found that most gram-positive bacteria (tainted) were sensitive to the dye, while most gram-negative bacteria (not tainted) were not, and observed that the dye tended to act as a bacteriostatic agent rather than a bactericide.
Precautions
One study in mice demonstrated dose-related carcinogenic potential at several different organ sites. The Food and Drug Administration in the US (FDA) has determined that gentian violet has not been shown by adequate scientific data to be safe for use in animal feed. Use of gentian violet in animal feed causes the feed to be adulterated and is a violation of the Federal Food, Drug, and Cosmetic Act in the US. On June 28, 2007, the FDA issued an "import alert" on farm raised seafood from China because unapproved antimicrobials, including gentian violet, had been consistently found in the products. The FDA report states:
"Like MG (malachite green), CV (crystal violet) is readily absorbed into fish tissue from water exposure and is reduced metabolically by fish to the leuco moiety, leucocrystal violet (LCV). Several studies by the National Toxicology Program reported the carcinogenic and mutagenic effects of crystal violet in rodents. The leuco form induces renal, hepatic and lung tumor in mice."
In 2019, Health Canada found medical devices that use gentian violet to be safe for use but recommended to stop using all drug products that contain gentian violet, including on animals, causing Canadian engineering schools to revisit the usage of this dye during orientation.
| Physical sciences | Organic salts | Chemistry |
1506087 | https://en.wikipedia.org/wiki/Epigenome | Epigenome | In biology, the epigenome of an organism is the collection of chemical changes to its DNA and histone proteins that affects when, where, and how the DNA is expressed; these changes can be passed down to an organism's offspring via transgenerational epigenetic inheritance. Changes to the epigenome can result in changes to the structure of chromatin and changes to the function of the genome. The human epigenome, including DNA methylation and histone modification, is maintained through cell division (both mitosis and meiosis). The epigenome is essential for normal development and cellular differentiation, enabling cells with the same genetic code to perform different functions. The human epigenome is dynamic and can be influenced by environmental factors such as diet, stress, and toxins.
The epigenome is involved in regulating gene expression, development, tissue differentiation, and suppression of transposable elements. Unlike the underlying genome, which remains largely static within an individual, the epigenome can be dynamically altered by environmental conditions.
Types
The main types of epigenetic changes include:
DNA methylation
Addition of a methyl group to the DNA molecule, typically at cytosine bases. This modification generally leads to gene silencing by preventing the binding of transcription factors and other proteins necessary for gene expression.
DNA functionally interacts with a variety of epigenetic marks, such as cytosine methylation, also known as 5-methylcytosine (5mC). This epigenetic mark is widely conserved and plays major roles in the regulation of gene expression, in the silencing of transposable elements and repeat sequences.
Individuals differ with their epigenetic profile, for example the variance in CpG methylation among individuals is about 42%. On the contrary, epigenetic profile (including methylation profile) of each individual is constant over the course of a year, reflecting the constancy of our phenotype and metabolic traits. Methylation profile, in particular, is quite stable in a 12-month period and appears to change more over decades.
Methylation sites
CoRSIVs are Correlated Regions of Systemic Interindividual Variation in DNA methylation. They span only 0.1% of the human genome, so they are very rare; they can be inter-correlated over long genomic distances (>50 kbp). CoRSIVs are also associated with genes involved in a lot of human disorders, including tumors, mental disorders and cardiovascular diseases. It has been observed that disease-associated CpG sites are 37% enriched in CoRSIVs compared to control regions and 53% enriched in CoRSIVs relative to tDMRs (tissue specific Differentially Methylated Regions).
Most of the CoRSIVs are only 200 – 300 bp long and include 5–10 CpG dinucleotides, the largest span several kb and involve hundreds of CpGs. These regions tend to occur in clusters and the two genomic areas of high CoRSIV density are observed at the major histocompatibility (MHC) locus on chromosome 6 and at the pericentromeric region on the long arm of chromosome 20.
CoRSIVs are enriched in intergenic and quiescent regions (e.g. subtelomeric regions) and contain many transposable elements, but few CpG islands (CGI) and transcription factor binding sites. CoRSIVs are under-represented in the proximity of genes, in heterochromatic regions, active promoters, and enhancers. They are also usually not present in highly conserved genomic regions.
CoRSIVs can have a useful application: measurements of CoRSIV methylation in one tissue can provide some information about epigenetic regulation in other tissues, indeed we can predict the expression of associated genes because systemic epigenetic variants are generally consistent in all tissues and cell types.
Factors affecting methylation pattern
Quantification of the heritable basis underlying population epigenomic variation is also important to delineate its cis- and trans-regulatory architecture. In particular, most studies state that inter-individual differences in DNA methylation are mainly determined by cis-regulatory sequence polymorphisms, probably involving mutations in TFBSs (Transcription Factor Binding Sites) with downstream consequences on local chromatin environment. The sparsity of trans-acting polymorphisms in humans suggests that such effects are highly deleterious. Indeed, trans-acting factors are expected to be caused by mutations in chromatin control genes or other highly pleiotropic regulators. If trans-acting variants do exist in human populations, they probably segregate as rare alleles or originate from somatic mutations and present with clinical phenotypes, as is the case in many cancers.
Correlation between methylation and gene expression
DNA methylation (in particular in CpG regions) is able to affect gene expression: hypermethylated regions tend to be differentially expressed. In fact, people with a similar methylation profile tend to also have the same transcriptome. Moreover, one key observation from human methylation is that most functionally relevant changes in CpG methylation occur in regulatory elements, such as enhancers.
Anyway, differential expression concerns only a slight number of methylated genes: only one fifth of genes with CpG methylation shows variable expression according to their methylation state. It is important to notice that methylation is not the only factor affecting gene regulation.
Methylation in embryos
It was revealed by immunostaining experiments that in human preimplantation embryos there is a global DNA demethylation process. After fertilisation, the DNA methylation level decreases sharply in the early pronuclei. This is a consequence of active DNA demethylation at this stage. But global demethylation is not an irreversible process, in fact de novo methylation occurring from the early to mid-pronuclear stage and from the 4-cell to the 8-cell stage.
The percentage of DNA methylation is different in oocytes and in sperm: the mature oocyte has an intermediate level of DNA methylation (72%), instead the sperm has high level of DNA methylation (86%). Demethylation in paternal genome occurs quickly after fertilisation, whereas the maternal genome is quite resistant at the demethylation process at this stage. Maternal different methylated regions (DMRs) are more resistant to the preimplantation demethylation wave.
CpG methylation is similar in germinal vesicle (GV) stage, intermediate metaphase I (MI) stage and mature metaphase II (MII) stage. Non-CpG methylation continues to accumulate in these stages.
Chromatin accessibility in germline was evaluated by different approaches, like scATAC-seq and sciATAC-seq, scCOOL-seq, scNOMe-seq and scDNase-seq. Stage-specific proximal and distal regions with accessible chromatin regions were identified. Global chromatin accessibility is found to gradually decrease from the zygote to the 8-cell stage and then increase. Parental allele-specific analysis shows that paternal genome becomes more open than the maternal genome from the late zygote stage to the 4-cell stage, which may reflect decondensation of the paternal genome with replacement of protamines by histones.
Sequence-Dependent Allele-Specific Methylation
DNA methylation imbalances between homologous chromosomes show sequence-dependent behavior. Difference in the methylation state of neighboring cytosines on the same chromosome occurs due to the difference in DNA sequence between the chromosomes. Whole-genome bisulfite sequencing (WGBS) is used to explore sequence-dependent allele-specific methylation (SD-ASM) at a single-chromosome resolution level and comprehensive whole-genome coverage. The results of WGBS tested on 49 methylomes revealed CpG methylation imbalances exceeding 30% differences in 5% of the loci.
On the sites of gene regulatory loci bound by transcription factors the random switching between methylated and unmethylated states of DNA was observed. This is also referred as stochastic switching and it is linked to selective buffering of gene regulatory circuit against mutations and genetic diseases. Only rare genetic variants show the stochastic type of gene regulation.
The study made by Onuchic et al. was aimed to construct the maps of allelic imbalances in DNA methylation, gene transcription, and also of histone modifications. 36 cell and tissue types from 13 participant donors were used to examine 71 epigenomes. The results of WGBS tested on 49 methylomes revealed CpG methylation imbalances exceeding 30% differences in 5% of the loci. The stochastic switching occurred at thousands of heterozygous regulatory loci that were bound to transcription factors. The intermediate methylation state is referred to the relative frequencies between methylated and unmethylated epialleles. The epiallele frequency variations are correlated with the allele affinity for transcription factors.
The analysis of the study suggests that human epigenome in average covers approximately 200 adverse SD-ASM variants. The sensitivity of the genes with tissue-specific expression patterns gives the opportunity for the evolutionary innovation in gene regulation.
Haplotype reconstruction strategy is used to trace chromatin chemical modifications (using ChIP-seq) in a variety of human tissues. Haplotype-resolved epigenomic maps can trace allelic biases in chromatin configuration. A substantial variation among different tissues and individuals is observed. This allows the deeper understanding of cis-regulatory relationships between genes and control sequences.
Histone modification
Post-translational modifications of histone proteins, which include methylation, acetylation, phosphorylation, ubiquitination, and sumoylation. These modifications can either activate or repress gene expression by altering chromatin structure and accessibility of the DNA to transcriptional machinery.
The epigenetic profiles of human tissues reveals the following distinct histone modifications in different functional areas:
Acetylation
Histone acetylation neutralizes the positive charge on histones. This weakens the electrostatic attraction to negatively charged DNA and causes unwinding of DNA from histones, making the DNA more accessible to the transcriptional machinery and hence resulting in transcriptional activation.
Methylation
Can lead to activation or repression of gene expression depending on the specific amino acids that are methylated.
Non-coding RNA gene silencing
Non-coding RNA (ncRNA) gene silencing involves various types of non-coding RNAs, such as microRNAs (miRNAs), long non-coding RNAs (lncRNAs), and small interfering RNAs (siRNAs). These RNA molecules can modulate gene expression by various mechanisms, including mRNA degradation, inhibition of translation, and chromatin remodeling.
Structural modifications
During the last few years, several methods have been developed to study the structural and consequently the functional modifications of chromatin. The first project that used epigenomic profiling to identify regulatory elements in the human genome was ENCODE (Encyclopedia of DNA Elements) that focused on profiling histone modifications on cell lines. A few years later ENCODE was included in the International Human Epigenome Consortium (IHEC), which aims to coordinate international epigenome studies.
The structural modifications that these projects aim to study can be divided into five main groups:
Nucleosome occupancy to detect regions with regulatory genes;
Chromatin interactions and domains;
Topological associated domains (TADs)
Topological associated domains are a degree of structural organization of the genome of the cell. They are formed by regions of chromatin, sized from 100 kilobases up to megabases, which highly self-interact. The domains are linked by other genomic regions, which, based on their size, are either called “topological boundary regions” or “unorganized chromatin”. These boundary regions separate the topological domains from heterochromatin, and prevent the amplification of the latter. Topological domains are diffused in mammalian, although similar genome partitions were identified also in Drosophila.
Topological domains in humans, like in other mammalians, have many functions regarding gene expression and transcriptional control process. Inside these domains, the chromatin shows to be well tangled, while in the boundary regions chromatin interactions are far less present. These boundary areas in particular show some peculiarity that determine the functions of all the topological domains.
Firstly, they contain insulator regions and barrier elements, both of which function as inhibitors of further transcription from the RNA polymerase enzyme. Such elements are characterized by the massive presence of insulator binding proteins CTCF.
Secondly, boundary regions block heterochromatin spreading, thus preventing the loss of useful genetic informations. This information derives from the observation that the heterochromatin mark H3K9me3 sequences clearly interrupts near boundary sequences.
Thirdly, transcription start sites (TSS), housekeeping genes and tRNA genes are particularly abundant in boundary regions, denoting that those areas have a prolific transcriptional activity, thanks to their structural characteristics, different from other topological regions.
Finally, in the border areas of the topological domains and their surroundings there is an enrichment of Alu/B1 and B2 SINE retrotransposons. In the recent years, those sequences were referred to alter binding site of CTCF, thus interfering with expression of some genomic areas.
Further proofs towards a role in genetic modulation and transcription regulation refers to the great conservation of the boundary pattern across mammalian evolution, with a dynamic range of small diversities inside different cell types, suggesting that these topological domains take part in cell-type specific regulatory events.
Correlation between methylation and 3D structure
The 4D Nucleome project aims to realize a 3D maps of mammalian genomes in order to develop predictive models to correlate epigenomic modifications with genetic variation. In particular the goal is to link genetic and epigenomic modifications with the enhancers and promoters which they interact with in three-dimensional space, thus discovering gene-set interactomes and pathways as new candidates for functional analysis and therapeutic targeting.
Hi-C is an experimental method used to map the connections between DNA fragments in three-dimensional space on a genome-wide scale. This technique combines chemical crosslinking of chromatin with restriction enzyme digestion and next-generation DNA sequencing.
This kind of studies are currently limited by the lack or unavailability of raw data.
Clinical significance
Cancer
Epigenetics is a currently active topic in cancer research. Human tumors undergo a major disruption of DNA methylation and histone modification patterns. The aberrant epigenetic landscape of the cancer cell is characterized by a global genomic hypomethylation, CpG island promoter hypermethylation of tumor suppressor genes, an altered histone code for critical genes and a global loss of monoacetylated and trimethylated histone H4.
Aging
The idea that DNA damage drives aging by compromising transcription and DNA replication has been widely supported since it was initially developed the 1980s. In recent decades, evidence has accumulated supporting the additional idea that DNA damage and repair elicit widespread epigenome alterations that also contribute to aging (e.g.). Such epigenome changes include age-related changes in the patterns of DNA methylation and histone modification.
Research
As a prelude to a potential Human Epigenome Project, the Human Epigenome Pilot Project aims to identify and catalogue Methylation Variable Positions (MVPs) in the human genome. Advances in sequencing technology now allow for assaying genome-wide epigenomic states by multiple molecular methodologies. Micro- and nanoscale devices have been constructed or proposed to investigate the epigenome.
An international effort to assay reference epigenomes commenced in 2010 in the form of the International Human Epigenome Consortium (IHEC). IHEC members aim to generate at least 1,000 reference (baseline) human epigenomes from different types of normal and disease-related human cell types.
Roadmap epigenomics project
One goal of the NIH Roadmap Epigenomics Project is to generate human reference epigenomes from normal, healthy individuals across a large variety of cell lines, primary cells, and primary tissues. Data produced by the project, which can be browsed and downloaded from the Human Epigenome Atlas, fall into five types that assay different aspects of the epigenome and outcomes of epigenomic states (such as gene expression):
Histone Modifications – Chromatin Immunoprecipitation Sequencing (ChIP-Seq) identifies genome wide patterns of histone modifications using antibodies against the modifications.
DNA Methylation – Whole Genome Bisulfite-Seq, Reduced Representation Bisulfite-Seq (RRBS), Methylated DNA Immunoprecipitation Sequencing (MeDIP-Seq), and Methylation-sensitive Restriction Enzyme Sequencing (MRE-Seq) identify DNA methylation across portions of the genome at varying levels of resolution down to basepair level.
Chromatin Accessibility – DNase I hypersensitive sites Sequencing (DNase-Seq) uses the DNase I enzyme to find open or accessible regions in the genome.
Gene Expression – RNA-Seq and expression arrays identify expression levels or protein coding genes.
Small RNA Expression – smRNA-Seq identifies expression of small noncoding RNA, primarily miRNAs.
Reference epigenomes for healthy individuals will enable the second goal of the Roadmap Epigenomics Project, which is to examine epigenomic differences that occur in disease states such as Alzheimer's disease.
| Biology and health sciences | Genetics | Biology |
1508434 | https://en.wikipedia.org/wiki/Coplanarity | Coplanarity | In geometry, a set of points in space are coplanar if there exists a geometric plane that contains them all. For example, three points are always coplanar, and if the points are distinct and non-collinear, the plane they determine is unique. However, a set of four or more distinct points will, in general, not lie in a single plane.
Two lines in three-dimensional space are coplanar if there is a plane that includes them both. This occurs if the lines are parallel, or if they intersect each other. Two lines that are not coplanar are called skew lines.
Distance geometry provides a solution technique for the problem of determining whether a set of points is coplanar, knowing only the distances between them.
Properties in three dimensions
In three-dimensional space, two linearly independent vectors with the same initial point determine a plane through that point. Their cross product is a normal vector to that plane, and any vector orthogonal to this cross product through the initial point will lie in the plane. This leads to the following coplanarity test using a scalar triple product:
Four distinct points, , are coplanar if and only if,
which is also equivalent to
If three vectors are coplanar, then if (i.e., and are orthogonal) then
where denotes the unit vector in the direction of . That is, the vector projections of on and on add to give the original .
Coplanarity of points in n dimensions whose coordinates are given
Since three or fewer points are always coplanar, the problem of determining when a set of points are coplanar is generally of interest only when there are at least four points involved. In the case that there are exactly four points, several ad hoc methods can be employed, but a general method that works for any number of points uses vector methods and the property that a plane is determined by two linearly independent vectors.
In an -dimensional space where , a set of points are coplanar if and only if the matrix of their relative differences, that is, the matrix whose columns (or rows) are the vectors is of rank 2 or less.
For example, given four points
if the matrix
is of rank 2 or less, the four points are coplanar.
In the special case of a plane that contains the origin, the property can be simplified in the following way:
A set of points and the origin are coplanar if and only if the matrix of the coordinates of the points is of rank 2 or less.
Geometric shapes
A skew polygon is a polygon whose vertices are not coplanar. Such a polygon must have at least four vertices; there are no skew triangles.
A polyhedron that has positive volume has vertices that are not all coplanar.
| Mathematics | Two-dimensional space | null |
1509289 | https://en.wikipedia.org/wiki/Magnetostatics | Magnetostatics | Magnetostatics is the study of magnetic fields in systems where the currents are steady (not changing with time). It is the magnetic analogue of electrostatics, where the charges are stationary. The magnetization need not be static; the equations of magnetostatics can be used to predict fast magnetic switching events that occur on time scales of nanoseconds or less. Magnetostatics is even a good approximation when the currents are not static – as long as the currents do not alternate rapidly. Magnetostatics is widely used in applications of micromagnetics such as models of magnetic storage devices as in computer memory.
Applications
Magnetostatics as a special case of Maxwell's equations
Starting from Maxwell's equations and assuming that charges are either fixed or move as a steady current , the equations separate into two equations for the electric field (see electrostatics) and two for the magnetic field. The fields are independent of time and each other. The magnetostatic equations, in both differential and integral forms, are shown in the table below.
Where ∇ with the dot denotes divergence, and B is the magnetic flux density, the first integral is over a surface with oriented surface element . Where ∇ with the cross denotes curl, J is the current density and is the magnetic field intensity, the second integral is a line integral around a closed loop with line element . The current going through the loop is .
The quality of this approximation may be guessed by comparing the above equations with the full version of Maxwell's equations and considering the importance of the terms that have been removed. Of particular significance is the comparison of the term against the term. If the term is substantially larger, then the smaller term may be ignored without significant loss of accuracy.
Re-introducing Faraday's law
A common technique is to solve a series of magnetostatic problems at incremental time steps and then use these solutions to approximate the term . Plugging this result into Faraday's Law finds a value for (which had previously been ignored). This method is not a true solution of Maxwell's equations but can provide a good approximation for slowly changing fields.
Solving for the magnetic field
Current sources
If all currents in a system are known (i.e., if a complete description of the current density is available) then the magnetic field can be determined, at a position r, from the currents by the Biot–Savart equation:
This technique works well for problems where the medium is a vacuum or air or some similar material with a relative permeability of 1. This includes air-core inductors and air-core transformers. One advantage of this technique is that, if a coil has a complex geometry, it can be divided into sections and the integral evaluated for each section. Since this equation is primarily used to solve linear problems, the contributions can be added. For a very difficult geometry, numerical integration may be used.
For problems where the dominant magnetic material is a highly permeable magnetic core with relatively small air gaps, a magnetic circuit approach is useful. When the air gaps are large in comparison to the magnetic circuit length, fringing becomes significant and usually requires a finite element calculation. The finite element calculation uses a modified form of the magnetostatic equations above in order to calculate magnetic potential. The value of can be found from the magnetic potential.
The magnetic field can be derived from the vector potential. Since the divergence of the magnetic flux density is always zero,
and the relation of the vector potential to current is:
Magnetization
Strongly magnetic materials (i.e., ferromagnetic, ferrimagnetic or paramagnetic) have a magnetization that is primarily due to electron spin. In such materials the magnetization must be explicitly included using the relation
Except in the case of conductors, electric currents can be ignored. Then Ampère's law is simply
This has the general solution
where is a scalar potential. Substituting this in Gauss's law gives
Thus, the divergence of the magnetization, has a role analogous to the electric charge in electrostatics and is often referred to as an effective charge density .
The vector potential method can also be employed with an effective current density
| Physical sciences | Magnetism | null |
1509614 | https://en.wikipedia.org/wiki/Black%20turtle%20bean | Black turtle bean | The black turtle bean is a small, shiny variety of the common bean (Phaseolus vulgaris) especially popular in Latin American cuisine, though it can also be found in the Cajun and Creole cuisines of south Louisiana. Like all varieties of the common bean, it is native to the Americas, but has been introduced around the world. It is also used in Indian cuisine, Tamil cuisine, where it is known as karuppu kaaramani and in Maharashtrian cuisine, where it is known as kala ghevada. It is widely used in Uttrakhand, where it is also known as "bhatt". It is a rich source of iron and protein. The black turtle bean is often simply called the black bean (, , , , , or in Spanish; and in Portuguese), although this terminology can cause confusion with at least three other types of black beans.
The black turtle bean is the only type of turtle bean. It is called turtle because of its hard outer "shell".
It is not to be confused with douchi, the Chinese dish made with black hulled soybeans.
Background
The black bean has a dense, meaty texture, which makes it popular in vegetarian dishes, such as frijoles negros and the Mexican-American black bean burrito. It is a very popular bean in various regions of Brazil, and is used in the national dish, feijoada. It is also a main ingredient of Moros y Cristianos in Cuba, is a required ingredient in the typical gallo pinto of Costa Rica and Nicaragua, is a fundamental part of Pabellón criollo in Venezuela, and is served in almost all of Latin America, as well as many Hispanic enclaves in the United States. In the Dominican Republic cuisine, it is also used for a variation of the Moros y Cristianos simply called Moro de habichuelas negras. The black turtle bean is also popular as a soup ingredient. In Cuba, black bean soup is a traditional dish, usually served with white rice. Black beans sticky rice is a Thai dessert.
The bean was first widely grown in the present-day United States after the Mexican–American War (1846–1848). However, initially the variety was primarily grown as a snap pea (for the edible seed pod).
It is also common to keep the boiled water of these beans (which acquires a black coloring) and consume it as a soup with other ingredients for seasoning (known as sopa negra, black soup, or as sopa de frijoles, bean soup), as a broth (caldo de frijol, bean broth) or to season or color other dishes (aforementioned gallo pinto, for example).
Samples of black turtle beans were reported in 2006 to contain total anthocyanins in their dried seed coats of 0−2.78 mg/g.
Gallery
| Biology and health sciences | Pulses | Plants |
1510249 | https://en.wikipedia.org/wiki/Bering%20Strait%20crossing | Bering Strait crossing | A Bering Strait crossing is a hypothetical bridge or tunnel that would span the relatively narrow and shallow Bering Strait between the Chukotka Peninsula in Russia and the Seward Peninsula in the U.S. state of Alaska. The crossing would provide a connection linking the Americas and Afro-Eurasia.
With the two Diomede Islands between the peninsulas, the Bering Strait could be spanned by a bridge or tunnel.
There have been several proposals for a Bering Strait crossing made by various individuals and media outlets. The names used for them include "The Intercontinental Peace Bridge" and "EurasiaAmerica Transport Link". Tunnel names have included "TKMWorld Link", "AmerAsian Peace Tunnel" and InterBering. In April 2007, Russian government officials told the press that the Russian government would back a US$65 billion plan by a consortium of companies to build a Bering Strait tunnel.
History
19th century
The concept of an overland connection crossing the Bering Strait goes back before the 20th century. William Gilpin, first governor of the Colorado Territory, envisaged a vast "Cosmopolitan Railway" in 1890 linking the entire world through a series of railways.
Two years later, Joseph Strauss, who went on to design over 400 bridges, and then serve as the project engineer for the Golden Gate Bridge, put forward the first proposal for a Bering Strait rail bridge in his senior thesis. The project was presented to the government of the Russian Empire, but it was rejected.
20th century
In 1904, a syndicate of American railroad magnates proposed (through a French spokesman) a SiberianAlaskan railroad from Cape Prince of Wales in Alaska through a tunnel under the Bering Strait and across northeastern Siberia to Irkutsk via Cape Dezhnev, Verkhnekolymsk, and Yakutsk (around of railroad to build, plus over in North America). The proposal was for a 90-year lease, and exclusive mineral rights for each side of the right-of-way. It was debated by officials and finally turned down on March 20, 1907.
Czar Nicholas II approved the American proposal in 1905 (only as a permission, not much financing from the Czar). Its cost was estimated at $65 million and $300 million, including all the railroads. These hopes were dashed with the outbreak of the 1905 Russian Revolution followed by World War I.
A Nazi plan to create a wide-gauge railroad called the Breitspurbahn was mooted to connect the cities of Europe, India, China and ultimately North America via the Bering Strait. The railroad was never built.
Interest was renewed during World War II with the completion in 19421943 of the Alaska Highway, linking the remote territory of Alaska with Canada and the continental United States. In 1942, the Foreign Policy Association envisioned the highway continuing to link with Nome near the Bering Strait, linked by highway to the railhead at Yakutsk, using an alternative sea-and-air ferry service across the Bering Strait. At the same time the road on the Russian side was extended by building the Kolyma Highway.
In 1958, engineer Tung-Yen Lin suggested the construction of a bridge across the Bering Strait "to foster commerce and understanding between the people of the United States and the Soviet Union". Ten years later he organized the Inter-Continental Peace Bridge, Inc., a non-profit institution organized to further this proposal. At that time he made a feasibility study of a Bering Strait bridge and estimated the cost to be $1 billion for the span. In 1994 he updated the cost to more than $4 billion. Like Gilpin, Lin envisioned the project as a symbol of international cooperation and unity, and dubbed the project the Intercontinental Peace Bridge.
21st century
According to a report in the Beijing Times in May 2014, Chinese transport experts had proposed building a roughly high-speed rail line from northeast China to the United States. The project would include a tunnel under the Bering Strait and connect to the contiguous United States via Wales, Alaska, along the river to Fairbanks, Alaska, and along the Alaska Highway to Edmonton, Alberta, Canada.
Several American entrepreneurs have also advanced private-sector proposals, such as an Alaska-based limited-liability company InterBering founded in 2010 to lobby for a cross-straits connection, and a 2018 cryptocurrency offering to fund the construction of a tunnel. In 2005, investor Neil Bush, younger brother of U.S. President George W. Bush and son of President George H. W. Bush, traveled abroad with Sun Myung Moon of the Unification Church as he promoted a proposal to dig a transportation corridor beneath the Bering Strait. When questioned by Mother Jones during the Republican primary campaign of his brother Jeb Bush a decade later in 2015, he denied having supported the tunnel project and said that he had traveled with Moon because he supported "efforts by faith leaders to call their flock into service to others."
Strategic military concerns
Proposals to build a crossing predate the Russian invasion of Ukraine and the Russian-Ukrainian War, which started in February 2022. It is not known how those events have affected strategic concerns relating to the proposed crossing, which would facilitate access by Russia to North America. Even before the invasion, commentators on the proposed link have flagged strategic military concerns as a factor in any decision to build the crossing.
Technical concerns
Distance
The straight distance between Russia and Alaska is . If building bridges and using the Diomede Islands, the straight distance over water for the three parts would be , and , in total .
Depth of water
The depth of the water is a minor problem, as the strait is no deeper than , comparable to the English Channel. The tides and currents in the area are not severe.
Weather-related challenges
Restrictions on construction work
The route is just south of the Arctic Circle, and the location has long, dark winters and extreme weather, including average winter lows of and temperatures approaching in cold snaps. This would mean that construction work would likely be restricted to five months of the year, around May to September, and centered during summer.
Exposed steel
The weather also poses challenges to exposed steel. In Lin's design, concrete covers all structures, to simplify maintenance and to offer additional stiffening.
Ice floes
Although there are no icebergs in the Bering Strait, ice floes up to thick are in constant motion during certain seasons, which could produce forces on the order of on a pier.
Tundra in surrounding regions
Roads on either side of the strait would likely have to cross tundra, requiring either an unpaved road or some way to avoid the effects of permafrost.
Likely route and expenses
Bridge option
If the crossing is chosen as a bridge, it would probably connect Wales, Alaska, to a location south of Uelen. The bridge would also likely be divided by the Diomede Islands, which are at the middle of the Bering Strait.
In 1994, Lin estimated the cost of a bridge to be "a few billion" dollars. The roads and railways on each side were estimated to cost $50 billion. Lin contrasted this cost to petroleum resources "worth trillions". Discovery Channel's Extreme Engineering estimates the cost of a highway, electrified double-track high-speed rail, and pipelines at $105 billion (in 2007 US dollars), five times the original cost of the 1994 Channel Tunnel.
Connections to the rest of the world
This excludes the cost of new roads and railways to reach the bridge. Aside from the technical challenges of building two bridges or a more than tunnel across the strait, another major challenge is that, , there is nothing on either side of the Bering Strait to connect the bridge to.
Russian side
The Russian side of the strait, in particular, is severely lacking in infrastructure. No railways exist for over in any direction from the strait.
The nearest major connecting highway is the M56 Kolyma Highway, which is currently unpaved and around from the strait. However, by 2042, the Anadyr Highway is expected to be completed connecting Ola and Anadyr, which is only about from the strait.
U.S. side
On the U.S. side, an estimated of highways or railroads would have to be built around Norton Sound, through a pass along the Unalakleet River, and along the Yukon River to connect to Manley Hot Springs Road – in other words, a route similar to that of the Iditarod Trail Race. A project to connect Nome, from the strait, to the rest of Alaska by a paved highway (part of Alaska Route 2) has been proposed by the Alaskan state government, although the very high cost ($2.3 to $2.7 billion, about $3 million per kilometer, or $5 million per mile) has so far prevented construction.
In 2016, the Alaskan road network was extended westwards by to Tanana, from the strait, by building a fairly simple road. The Alaska Department of Transportation & Public Facilities project was supported by local indigenous groups such as the Tanana Tribal Council.
Track gauge
Another complicating factor is the different track gauges in use. Mainline rail in the US, Canada, China, and the Koreas uses standard gauge of 1435 millimeters. Russia uses the slightly broader Russian gauge of 1520 mm. Solutions to this break of gauge include:
To have all cargo in containers, which are fairly easily reloaded from one train to another. This is used on the increasingly popular China–Europe rail freight route, which has two breaks of gauge. It is possible to transfer a 60-container train in one hour.
Another solution is variable gauge axles for locomotives and rolling stock, such as those made by Talgo. A gauge changer modifies the gauge of the wheels while the train traverses the GC equipment at a speed of , which is about 4 seconds per railcar. This is faster than is possible with the transfer of ISO containers.
The TKMWorld Link
The TKMWorld Link (Russian: ТрансКонтинентальная магистраль, English: Transcontinental Railway), also called ICL-World Link (Intercontinental link), was a planned link between Siberia and Alaska to deliver oil, natural gas, electricity, and rail passengers to the United States from Russia. Proposed in 2007, the plan included provisions to build a tunnel under the Bering Strait, which, if built, would have been the longest tunnel in the world, surpassing the Line 3 (Guangzhou Metro) tunnel. The tunnel was intended to be part of a railway joining Yakutsk, the capital of the Russian republic of Yakutia, and Komsomolsk-on-Amur, in the Russian Far East, with the western coast of Alaska. The Bering Strait tunnel was estimated to cost between $10 billion and $12 billion, while the entire project was estimated to cost $65 billion.
In 2008, Russian Prime Minister Vladimir Putin approved the plan to build a railway to the Bering Strait area, as a part of the development plan to run until 2030. The more than tunnel would have run under the Bering Strait between Chukotka, in the Russian far east, and Alaska. The cost was estimated as $66 billion.
In late August 2011, at a conference in Yakutsk in eastern Russia, the plan was backed by some of President Dmitry Medvedev's top officials, including Aleksandr Levinthal, the deputy federal representative for the Russian Far East. Supporters of the idea believed that it would be a faster, safer, and cheaper way to move freight around the world than container ships. They estimated it could carry about 3% of global freight and make about $7 billion a year. Shortly after, the Russian government approved the construction of the $65 billion Siberia-Alaska rail and tunnel across the Bering Strait.
Observers doubted that the rail link would be cheaper than ship, bearing in mind that the cost for rail transport from China to Europe is higher than by ship (except for expensive cargo where lead time is important).
In 2013, the Amur Yakutsk Mainline connecting the Yakutsk railway ( from the strait) with the Trans-Siberian Railway was completed. However, this railway is meant for freight and is too tightly curved for high-speed passenger trains. Future projects include the and Kolyma–Anadyr highway. The Kolyma–Anadyr highway has started construction, but will be a narrow gravel road.
USCanadaRussiaChina railway
In 2014, China was considering construction of a US-Canada-Russia-China bullet train that would include a undersea tunnel crossing the Bering Strait and would allow passengers to travel between the United States and China in about two days.
Although the press was skeptical of the project, China's state-run China Daily claimed that China possessed the necessary technology. It was unknown who was expected to pay for the construction, although China had in other projects offered to build and finance them, and expected the money back in the end through fees or rents.
Trans-Eurasian Belt Development
In 2015, another possible collaboration between China and Russia was reported, part of the Trans-Eurasian Belt Development, a transportation corridor across Siberia that would also include a road bridge with gas and oil pipelines between the easternmost point of Siberia and the westernmost point of Alaska. It would link London and New York by rail and superhighway via Russia if it were to go ahead.
China's Belt and Road Initiative has similar plans, so the project would work in parallel for both countries.
| Technology | Multi-modal crossings | null |
1511050 | https://en.wikipedia.org/wiki/Picea%20glauca | Picea glauca | Picea glauca (Moench) Voss., the white spruce, is a species of spruce native to the northern temperate and boreal forests in Canada and United States, North America.
Picea glauca is native from central Alaska all through the east, across western and southern/central Canada to the Avalon Peninsula in Newfoundland, Quebec, Ontario and south to Montana, North Dakota, Minnesota, Wisconsin, Michigan, Upstate New York and Vermont, along with the mountainous and immediate coastal portions of New Hampshire and Maine, where temperatures are just barely cool and moist enough to support it. There is also an isolated population in the Black Hills of South Dakota and Wyoming. It is also known as Canadian spruce, skunk spruce, cat spruce, Black Hills spruce, western white spruce, Alberta white spruce, and Porsild spruce.
Description
The white spruce is a large evergreen conifer which normally grows to tall, but can grow up to tall with a trunk diameter of up to . The bark is thin and scaly, flaking off in small circular plates across. The crown is narrowconical in young trees, becoming cylindrical in older trees. The shoots are pale buff-brown, glabrous in the east of the range, but often pubescent in the west, and with prominent pulvini. The leaves are needle-like, long, rhombic in cross-section, glaucous blue-green above (hence glauca) with several thin lines of stomata, and blue-white below with two broad bands of stomata.
The cones are pendulous, slender, cylindrical, long and 1.5 cm wide when closed, opening to 2.5 cm broad. They have thin, flexible scales 15 mm long with a smoothly rounded margin. They are green or reddish, maturing to pale brown 4 to 8 months after pollination. The seeds are black, 2 to 3 mm long, with a slender, tan wing 5 to 8 mm long.
Seeds
Seeds are small, 2.5 to 5 mm long, oblong, and acute at the base. Determinations of the average number of sound seeds per white spruce cone have ranged from 32 to 130.
Common causes of empty seed are lack of pollination, abortion of the ovule, and insect damage.
The average weight per individual seed varies from 1.1 mg to 3.2 mg.
Each seed is clasped by a thin wing 2 to 4 times as long as the seed. Seed and wing are appressed to the cone scale. Embryo and megagametophyte are soft and translucent at first; later the endosperm becomes firm and milky white, while the embryo becomes cream-coloured or light yellow. At maturity, the testa darkens rapidly from light brown to dark brown or black. Mature seeds "snaps in two" when cut by a sharp knife on a firm surface.
White spruce cones reach their maximum size after 800 GDD. Cone moisture content decreases gradually after about 1000 GDD.
Cone colour also can be used to help determine the degree of maturation, but cones may be red, pink or green. Collection and storage dates and conditions influence germination requirements and early seedling growth.
A bushel (35 L) of cones, which may contain 6,500 to 8,000 cones, yields of clean seed.
Seed dispersal begins after cone scales reflex with cone maturation in the late summer or early fall of the year of formation. Cones open at moisture contents of 45% to 70% and specific gravities of 0.6 to 0.8. Weather affects both the initiation and pattern of seed dispersal, but cone opening and the pattern of seed dispersal can vary among trees in the same stand. Even after dispersal has begun, cold, damp weather will cause cone scales to close; they will reopen during dry weather. Most seed falls early rather than late, but dispersal may continue through fall and winter and even into the next growing season. Seed dispersal occurs mainly in late summer or early fall.
White spruce seed is initially dispersed through the air by wind. Both the initiation and pattern of seed dispersal depend on the weather, but these can vary among trees in the same stand. Small amounts of white spruce seed are normally dispersed beyond 100 m from the seed source, but exceptionally seeds have been found more than 300–400 m from the nearest seed source.
Root system
The root system of white spruce is highly variable and adaptable, responding to a variety of edaphic factors, especially soil moisture, soil fertility, and mechanical impedance. On soils that limit rooting depth, the root system is plate-like, but it is a common misconception to assume that white spruce is genetically constrained to develop plate-like root systems irrespective of soil conditions. In the nursery, or naturally in the forest, white spruce usually develops several long 'running' roots just below the ground surface.
The structure of the tracheids in the long lateral roots of white spruce varies with soil nitrogen availability.
Stem
White spruce can live for several hundred years, with an estimated average lifespan of 250 to 300 years.
Slow-growing trees in rigorous climates are also capable of great longevity. White spruce high on the shore of Urquhart Lake, Northwest Territories, were found to be more than 300 years old.
Bark
The bark of mature white spruce is scaly or flaky, grey-brown or ash-brown, but silvery when freshly exposed. Resin blisters are normally lacking, but the Porsild spruce Picea glauca var. porsildii Raup has been credited with having smooth resin-blistered bark.
White spruce bark is mostly less than 8 mm and not more than 9.5 mm thick.
Chemistry
Isorhapontin can be found in spruce species such as the white spruce.
P. glauca has three different genomes; a nuclear genome, a mitochondrial genome, and a plastid (i.e. chloroplast) genome. The large (20 Gbp) nuclear genome of P. glauca (genotype WS77111) was published in 2015, and the organellar (plastid and mitochondrial) genomes (genotype PG29) were published in SD Jackman et al. 2015. The plastid genome of P. glauca (genotype WS77111) has also been published.
Varieties
Several geographical varieties have been described, but are not accepted as distinct by all authors. These comprise, from east to west:
Picea glauca var. glauca (typical or eastern white spruce): from Newfoundland west to eastern Alberta, on lowland plains.
Picea glauca var. densata (Black Hills white spruce): The Black Hills in South Dakota.
Picea glauca var. albertiana (Alberta white spruce): The Rocky Mountains in Alberta, British Columbia and northwest Montana.
Picea glauca var. porsildii (Alaska white spruce): Alaska and Yukon.
The two western varieties are distinguished by pubescent shoots, and may be related to extensive hybridisation and intergradation with the closely related Engelmann spruce found further south in the Rocky Mountains. White spruce also hybridises readily with the closely related Sitka spruce where they meet in southern Alaska and northwestern British Columbia; this hybrid is known as Picea × lutzii.
Distribution and habitat
White spruce has a transcontinental range in North America. In Canada, its contiguous distribution encompasses virtually the whole of the Boreal, Subalpine, Montane, Columbia, Great Lakes–St. Lawrence, and Acadian Forest Regions, extending into every province and territory. On the west coast of Hudson Bay, it extends to Seal River, about 59°N, "from which the northward limit runs apparently almost directly north-west to near the mouth of the Mackenzie River, or about latitude 68°". Collins and Sumner reported finding white spruce within 13 km of the Arctic coast in the Firth Valley, Yukon, at about 69°30′ N, 139°30′ W. It reaches within 100 km of the Pacific Ocean in the Skeena Valley, overlapping with the range of Sitka spruce (Picea sitchensis), and almost reaching the Arctic Ocean at latitude 69° N in the District of Mackenzie, with white spruce up to 15 m high occurring on some of the islands in the Delta near Inuvik. The wide variety of ecological conditions in which 4 Quebec conifers, including white spruce, are able to establish themselves, was noted by Lafond, but white spruce was more exacting than black spruce. In the United States, the range of white spruce extends into Maine, Vermont, New Hampshire, New York, Michigan, Wisconsin, Minnesota, and Alaska, where it reaches the Bering Strait in 66°44′ N" at Norton Bay and the Gulf of Alaska at Cook Inlet.
Southern outliers have been reported in southern Saskatchewan and the Cypress Hills of southwestern Saskatchewan and southeastern Alberta, northwestern Montana, south-central Montana, in the Black Hills on the Wyoming–South Dakota boundary, on the Manitoba–North Dakota boundary, and at Shushan, New York.
White spruce is the northernmost tree species in North America, reaching just north of 69°N latitude in the Mackenzie River delta. It grows between sea level and an elevation of . Its northern distribution roughly correlates to the location of the tree line, which includes an isothermic value of for mean temperature in July, as well as the position of the Arctic front; cumulative summer degree days, mean net radiation, and the amount of light intensities also figure. White spruce is generally found in regions where the growing season exceeds 60 days annually.
The southern edge of the zone in which white spruce forms 60% or more of the total stand corresponds more or less to the July isotherm of around the Great Lakes; in the Prairie Provinces its limit is north of this isotherm. During the summer solstice, day length values range from 17 hours at its southern limits to 24 hours above the Arctic Circle.
One of the hardiest conifers, white spruce in parts of its range withstands mean daily January temperature of and extreme minimum temperatures as low as ; minimum temperatures of are general throughout much of the range except the southernmost and southeasternmost parts. By itself, or with black spruce and tamarack (Larix laricina), white spruce forms the northern boundary of tree-form growth. White spruce up to 15 m in height occur at 69°N on islands in the Mackenzie Delta near Inuvik in the Northwest Territories. Hustich (1966) depicted Picea spp. as forming the northernmost limit of tree growth in North America.
The arctic or northern timberline in North America forms a broad transition zone from Labrador to northern Alaska. In Labrador, white spruce is not abundant and constitutes less than 5% of the forest, with a range that coincides very closely with that of black spruce but extending slightly further north.
The range of white spruce extends westwards from Newfoundland and Labrador, and along the northern limit of trees to Hudson Bay, Northwest Territories, Yukon, and into northwestern Alaska. Across western Canada and Alaska, white spruce occurs further north than black spruce, and, while poplar (Populus), willow, and birch may occur along streams well into the tundra beyond the limits of spruce, the hardwoods are usually no more than scrub. Spruce characteristically occurs in fingers of tree-form forest, extending far down the northern rivers and as scattered clumps of dwarfed "bush" spruce on intervening lands. In Manitoba, Scoggan noted that the northernmost collection of white spruce was at latitude 59°48’N, but Bryson et al. found white spruce in the northern edge of continuous forest in central Canada at Ennadai Lake, about 60°45′ N, 101°’W, just north of the northwest corner of Manitoba. Bryson et al. noted that the forest retained "the same general characteristics as when it was first described [by Tyrrell] in 1896". Collins and Sumner reported finding white spruce within 13 km of the Arctic coast in the Firth valley, Yukon, at about 69°30′ N, 139°30′ W, and Sargent noted that white spruce in Alaska "reached Behring Strait in 66°44′ N".
Climate, especially temperature, is obviously a factor in determining distributions of northern flora. Halliday and Brown suggested that white spruce's northern limit corresponds "very closely" with the July mean monthly isotherm of 10 °C in Ungava, but that the northern limit west of Hudson Bay was south of that isotherm. Other climatic factors that have been suggested as affecting the northern limit of white spruce include: cumulative summer degree days, position of the Arctic front in July, mean net radiation especially during the growing season, and low light intensities. Topography, soil conditions, and glaciation may also be important in controlling northern limits of spruce.
The southern limit of distribution of white spruce is more complex. From east of the main range of coastal mountains in British Columbia, the southern continuous limit of white spruce is the forest/prairie interface through Alberta, Saskatchewan, Manitoba, the northern parts of Minnesota and Wisconsin, central Michigan, northeastern New York, and Maine. Sargent and Harlow and Harrar also included Vermont and New Hampshire; and, while Dame and Brooks excluded New York and states further west, they included Massachusetts as far south as Amherst and Northampton, "probably the southern limit of the species" in that area. Nisbet gave the range of white spruce as extending to "Carolina", but he did not recognize red spruce as a species and presumably included it with white spruce.
Towards the southern parts of its range, white spruce encounters increasingly effective ecological competition from hardwoods, some of which may reinforce their growth-rate or sprouting competitiveness with allelopathic depredation of coniferous regeneration. Further southward extension of the distribution is inhibited by white spruce's cold requirement.
As an exotic species
As an exotic, white spruce is widespread but uncommon. It was introduced into England and parts of continental Europe in or soon after the year 1700, into Denmark about 1790, and into Tasmania and Ceylon shortly before 1932.
Nisbet noted that firmly-rooted white spruce served very well to stabilize windswept edges of woods in Germany. In a narrow belt of mixed Norway and white spruces over an extremely exposed hilltop crest at high elevation in northern England, the Norway spruce were "completely dwarfed" whereas the white spruce had reached heights of between 3 and 4.3 m. The age of the belt was not recorded, but adjoining 66-year-old stands may have been of the same vintage.
White spruce has also been used as a minor plantation species in England and Scotland. In Scotland, at Corrour, Inverness-shire, Sir John Stirling Maxwell in 1907 began using white spruce in his pioneering plantations at high elevations on deep peat. However, plantations in Britain have generally been unsatisfactory, mainly because of damage by spring frosts after mild weather had induced flushing earlier in the season. However, the species is held in high regard in the Belgian peat region, where it grows better than other spruces.
Ecology
White spruce is a climax canopy tree in the boreal forests of Canada and Alaska. It generally occurs on well-drained soils in alluvial and riparian zones, although it also occurs in soils of glacial and lacustrine origin. The understory is dominated by feather mosses (Hylocomium splendens, Pleurozium schreberi, Ptilium crista-castrensis) and fork mosses, and occasionally peat moss. In the far north, the total depth of the moss and underlying humus is normally between , although it tends to be shallower when hardwoods are present in the stand.
White spruce grows in soils with pH values of 4.7–7.0, although they have been found in soils as acidic as 4.0 in subalpine fir forests in the Northwest Territories. A presence of calcium in the soil is common to white spruce found in northern New York. White spruce most commonly grows in the soil orders of Alfisols and Inceptisols. Soil properties such as fertility, temperature, and structural stability are partial determinants of the ability of white spruce to grow in the extreme northern latitudes. In the northern limits of its range, white spruce is the climax species along with black spruce; birch and aspen are the early succession species.
Wildfires typically occur every 60 to 200 years, although they have been known to occur as infrequently as every 300 years.
White spruce will grow in USDA Growing Zones 3–7, but is not adapted to heat and humidity and will perform poorly in a hot climate. The tree attains its greatest longevity and growth potential in Zones 3–4.
Wildlife such as deer, rabbits, and grouse browse the foliage during the winter. The seeds are eaten by small mammals like the red squirrel and birds such as chickadee, nuthatch, and pine siskin.
Soils
White spruce occurs on a wide variety of soils, including soils of glacial, lacustrine, marine, and alluvial origins; overlying basic dolomites, limestones and acidic Precambrian and Devonian granites and gneisses; and Silurian sedimentary schists, shales, slates, and conglomerates. The wide range of textures accommodated includes clays, even those that are massive when wet and columnar when dry, sand flats, and coarse soils. Its occurrence on some organic soils is not characteristic, except perhaps on shallow mesic organic soils in Saskatchewan and in association with black spruce on organic soils in central Yukon.
Podzolized, brunisolic, luvisolic, gleysolic, and regosolic (immature) soils are typical of those supporting white spruce throughout the range of the species. Soils supporting white spruce are most commonly Alfisols or Inceptisols. In the podzol region of Wisconsin, white spruce occurs on loam podzols, podzolized gley loams, strongly podzolized clays, gley-podzol clays, stream-bottom soils, and wood peat. Moist sandy loams also support good growth. On sandy podzols, it is usually a minor species. Good development occurs on moist alluvium on the banks of streams and borders of swamps. White spruce makes good growth on well-drained lacustrine soils in Alberta Mixedwoods, on moderately-well-drained clay loams in Saskatchewan,, and on melanized loams and clays (with sparse litter and a dark-coloured organically-enriched mineral horizon) in the Algoma district of Ontario.
White spruce becomes less accommodating of soil with increasing severity of climate. The distribution of white spruce in Labrador seems to depend almost entirely on the character of the soil, and between the southwestern shores of Hudson Bay and the northeastern regions of Saskatchewan, white spruce is confined to very local physiographic features, characterized by well-drained or fertile soils.
On dry, deep, outwash deposits in northern Ontario, both white spruce and aspen grow slowly. But, broadly, white spruce is able to tolerate considerable droughtiness of sites that are fertile, and no fertile site is too moist unless soil moisture is stagnant. Soil fertility holds the key not just to white spruce growth but to the distribution of the species. At least moderate fertility is needed for good growth, but white spruce occurs on many sites where nutrient deficiencies depress its growth more than that of black spruce, red spruce, Norway spruce, and the pines generally. Minimum soil-fertility standards recommended for white spruce sufficient to produce 126 to 157 m3/ha of wood at 40 years are much higher than for pine species commonly planted in the Lake States (Wilde 1966): 3.5% organic matter, 12.0 meq/100 g exchange capacity, 0.12% total N, 44.8 kg/ha available P, 145.7 kg/ha available K, 3.00 meq/100 g exchangeable Ca, and 0.70 meq/100 g exchangeable Mg.
Forest floors under stands dominated by white spruce respond in ways that vary with site conditions, including the disturbance history of the site. Composition, biomass, and mineral soil physical and chemical properties are affected. In Alaska, the accumulation of organic layers (to greater thicknesses in mature stands of spruce than those in hardwood stands on similar sites) leads to decreased soil temperatures, in some cases leading to the development of permafrost. Acidity of the mineral soil sampled at an average depth of 17 cm in 13 white spruce stands on abandoned farmland in Ontario increased by 1.2 pH units over a period of 46 years.
A considerable range of soil pH is tolerated by white spruce. Thrifty stands of white spruce in Manitoba have developed on soils of pH 7.6 at only 10 cm below the surface, and pH 8.4 at 43 cm below the surface; rooting depth in those soils was at least 81 cm. An abundant calcium supply is common to most white spruce locations in New York state. Chlorosis was observed in young white spruce in heavily limed nursery soils at about pH 8.3. Wilde gave 4.7 to 6.5 as the approximate optimum range of pH for white spruce in Wisconsin, but optimum growth seems possible at pH levels up to 7.0 and perhaps higher. Alluvium on the floodplains of northern rivers shows pH levels from 5.0 to 8.2. High-lime ecotypes may exist, and in Canada Forest Section B8 the presence of balsam poplar and white spruce on some of the moulded moraines and clays seems to be correlated with the considerable lime content of these materials, while calcareous soils are favourable sites for northern outliers of white spruce.
Mature stands of white spruce in boreal regions often have well-developed moss layers dominated by feather mosses, e.g., Hylocomium splendens, Pleurozium schreberi, Ptlium crista-castrensis, and Dicranum, rather than Sphagnum. The thickness of the moss–organic layer commonly exceeds 25 cm in the far north and may approach twice that figure. The mosses compete for nutrients and have a major influence on soil temperatures in the rooting zone. Permafrost development in parts of Alaska, Yukon, and the Northwest Territories is facilitated by the insulative organic layer (Viereck 1970a, b, Gill 1975, Van Cleve and Yarie 1986).
Cold hardiness
White spruce is extremely hardy to low temperatures, provided the plant is in a state of winter dormancy. Throughout the greater part of its range, white spruce routinely survives and is undamaged by winter temperatures of , and even lower temperatures occur in parts of the range. Boreal Picea are among the few extremely hardy conifers in which the bud primordia are able to survive temperatures down to .
Especially important in determining the response of white spruce to low temperatures is the physiological state of the various tissues, notably the degree of "hardening" or dormancy. A natural progression of hardening and dehardening occurs in concert with the seasons. While different tissues vary in ability to tolerate exposure to stressful temperatures, white spruce, as with woody plants in general, has necessarily developed sufficient winter hardiness in its various tissues to enable them to survive the minimum temperatures experienced in the distribution range.
White spruce is subject to severe damage from spring frosts. Newly flushed shoots of white spruce are very sensitive to spring frost. This sensitivity is a major constraint affecting young trees planted without overstorey nurses in boreal climates.
Forest succession
Forest succession in its traditional sense implies two important features that resist direct examination. First, classical definitions generally connote directional changes in species composition and community structure through time, yet the time frame needed for documentation of change far exceeds an average human lifespan. The second feature that defies quantitative description is the end point or climax.
Floodplain deposits in the Northwest Territory, Canada, are important in relation to the development of productive forest types with a component of white spruce. The most recently exposed surfaces are occupied by sandbar vegetation or riparian shrub willows and Alnus incana. With increasing elevation, the shrubs give way successively to balsam poplar and white spruce forest. In contrast, older floodplains, with predominantly brown wooded soils, typically carry white spruce–trembling aspen mixedwood forest.
Interrelationships among nutrient cycling, regeneration, and subsequent forest development on floodplains in interior Alaska were addressed by Van Cleve et al., who pointed out that the various stages in primary succession reflect physical, chemical, and biological controls of ecosystem structure and function. Thus, each successional stage has a species combination in harmony with site quality. Short-circuiting succession by planting a late successional species such as white spruce on an early successional surface may result in markedly reduced growth rates because of nitrogen insufficiency. Without application of substantial amounts of fertilizer, use would have to be made of early successional alder and its site-ameliorating additions of nitrogen.
Neiland and Viereck noted that “the slow establishment and growth of spruce under birch stands [in Alaska] may be partially due to effects of shading and general competition for water and nutrients, but may also be more directly related to the birch itself. Heikinheimo found that birch ash inhibited white spruce seedlings, and Gregory found that birch litter has a smothering effect on spruce seedlings.".
On dry upland sites, especially south-facing slopes, the mature vegetation is white spruce, white birch, trembling aspen, or a combination of these species. Succession follows in one of two general patterns. In most cases, aspen and birch develop as a successional stage after fire before reaching the spruce stage. But, occasionally, with optimal site conditions and a source of seed, white spruce will invade with the hardwoods or within a few years thereafter, thereby producing even-aged white spruce stands without an intervening hardwood stage.
Associated forest cover
The White Spruce Cover Type may include other species in small numbers. In Alaska, associates include paper birch, trembling aspen, balsam poplar, and black spruce; in western Canada, additional associates are subalpine fir, balsam fir, Douglas-fir, jack pine, and lodgepole pine. Seral species giving way to white spruce include paper birch, aspen, balsam poplar, jack pine, and lodgepole pine. On certain river bottom sites, however, black spruce may replace white spruce. Earlier successional stages leading to the white spruce climax are the white spruce–paper birch, white spruce–aspen, balsam poplar, jack pine, and lodgepole pine types. The type shows little variation. The forest is generally closed and the trees well formed, other than those close to the timberline. Lesser vegetation in mature stands is dominated by mosses. Vascular plants are typically few, but shrubs and herbs that occur “with a degree of regularity” include: alder, willows, mountain cranberry, red-fruit bearberry, black crowberry, prickly rose, currant, buffaloberry, blueberry species, bunchberry, twinflower, tall lungwort, northern comandra, horsetail, bluejoint grass, sedge species, as well as ground-dwelling mosses and lichens. Several white spruce communities have been identified in interior Alaska: white spruce/feathermoss; white spruce/dwarf birch/feathermoss; white spruce/dwarf birch/sphagnum; white spruce/avens/moss; and white spruce/alder/bluejoint.
Of the Eastern Forest Cover Types recognized by the Society of American Foresters, only one, White Spruce, names that species in its title. The eastern White Spruce Cover Type, as defined, encompasses white spruce both in pure stands, and in mixed stands "in which white spruce is the major [undefined] component."
In most of its range, white spruce occurs more typically in association with trees of other species than in pure stands.
White spruce is an associated species in the following Eastern Forest cover types, by the Society of American Foresters; in the Boreal Forest Region: (1) jack pine, (5) balsam fir, (12) black spruce, (16) aspen, (18) paper birch, and (38) tamarack; in the Northern Forest Region: (15) red pine, (21) eastern white pine, (24) hemlock-yellow birch, (25) sugar maple-beech-yellow birch, (27) sugar maple, (30) red spruce-yellow birch, (32) red spruce, (33) red spruce-balsam fir, (37) northern white-cedar, and (39) black ash-American elm-red maple.
Predators
Outbreaks of spruce beetles have destroyed over of forests in Alaska.
Although sometimes described as relatively resistant to attack by insects and disease, white spruce is far from immune to depredation. Important insect pests of white spruce include the spruce budworm (Choristoneura fumiferana), the yellow-headed spruce sawfly (Pikonema alaskensis), the European spruce sawfly (Gilpinia hercyniae), the spruce bud moth (Zeiraphera canadensis), and spruce beetle (Dendroctonus rufipennis). As well, other budworms, sawflies, and bark beetles, gall formers, bud midges, leaf miners, aphids, leaf eaters, leaf rollers, loopers, mites, scales, weevils, borers, pitch moths, and spittlebugs cause varying degrees of damage to white spruce.
A number of sawflies feed on spruce trees. Among them European spruce sawfly, yellow-headed spruce sawfly, green-headed spruce sawfly and the spruce webspinning sawfly.
More than a dozen kinds of looper feed on the spruces, fir, and hemlock in eastern Canada. The full-grown larvae of the larvae vary in length from 15 mm to 35 mm. Some feed briefly in the fall and complete their feeding in the spring; others feed mainly in the summer; still others feed mainly in the late summer and fall.
The fall and spring feeding group includes the dash-lined looper (Protoboarmia porcelaria indicataria), the diamond-backed looper (Hypagyrtis piniata), the fringed looper (Campaea perlata), and the false loopers (Syngrapha species). The summer feeding group includes the false hemlock looper (Nepytia canosaria Walker), occasionally occurring in large numbers and usually in conjunction with the hemlock looper (Lambdina fiscellaria), the Eupithecia species, the yellowlined conifer looper (Cladara limitaria), and the saddleback looper (Ectropis crepuscularia).
The late summer and fall group includes the common spruce-fir looper (Semiothisa signaria dispuncta) and the similar hemlock angle (moth) Macaria fissinotata on hemlock, the small spruce loopers Eupithecia species, the gray spruce looper Caripeta divisata, occasionally abundant, the black-dashed hydriomena moth (Hydriomena divisaria), and the whitelined looper (Eufidonia notataria).
Cultivars
Numerous cultivars of various sizes, colours and shapes have been selected for use in parks and gardens. The following have gained the Royal Horticultural Society's Award of Garden Merit.
Picea glauca 'Echiniformis'
Picea glauca var. albertiana 'Alberta Globe'
Picea glauca var. albertiana 'Conica'
'Conica' is a dwarf conifer with very slender leaves, like those normally found only on one-year-old seedlings, and very slow growth, typically only per year. Older specimens commonly 'revert', developing normal adult foliage and starting to grow much faster; this 'reverted' growth must be pruned if the plant is to be kept dwarf.
Uses
The wood of white spruce is of a lower quality than that of Engelmann spruce, but is stronger. It was used to make shelters and as firewood by Native Americans and European settlers in Alaska, where lodgepole pine does not grow. The wood is of major economic importance in Canada, being harvested for paper and construction. It is also used as a Christmas tree.
The wood is also exported to Japan where, known as "shin-kaya", it is used to make go boards as a substitute for the rare kaya wood. Additionally, Picea glauca var. densata is used for bonsai.
White spruce is the provincial tree of Manitoba and the state tree of South Dakota.
The new growth or tips of white spruce is used in beer making, gin production, flavouring soda, candy making or in pickles and preserves.
| Biology and health sciences | Pinaceae | Plants |
1511315 | https://en.wikipedia.org/wiki/Picea%20engelmannii | Picea engelmannii | Picea engelmannii, with the common names Engelmann spruce, white spruce, mountain spruce, and silver spruce, is a species of spruce native to western North America. Highly prized for producing distinctive tone wood for acoustic guitars and other instruments, it is mostly a high-elevation mountain tree but also appears in watered canyons.
Description
Picea engelmannii is a medium-sized to large evergreen tree growing to tall, exceptionally to tall, and with a trunk diameter of up to . The reddish bark is thin and scaly, flaking off in small circular plates across. The crown is narrow conic in young trees, becoming cylindric in older trees. The shoots are buff-brown to orange-brown, usually densely pubescent, and with prominent pulvini. The leaves are needle-like, long, flexible, rhombic in cross-section, glaucous blue-green above with several thin lines of stomata, and blue-white below with two broad bands of stomata. The needles have a pungent odour when crushed.
Purple cones of about appear in spring, releasing yellow pollen when windy. The cones are pendulous, slender cylindrical, long and broad when closed, opening to broad. They have thin, flexible scales long, with a wavy margin. They are reddish to dark purple, maturing to light brown 4–7 months after pollination. The seeds are black, long, with a slender, long light brown wing.
The tree grows in a krummholz form along the fringe of alpine tundras.
Distribution
Engelmann spruce is mostly a higher-elevation mountain tree, in many areas reaching the tree line, but at lower elevations occupies cool watered canyons. It grows from above sea level, rarely lower towards the northwest.
Englemann spruce is native to western North America, common in the Rocky Mountains and east slopes of the Cascade Range from central British Columbia to Southern Oregon in the Cascades and Montana, Idaho, and Colorado, and more sparsely towards Arizona and New Mexico in the Sky islands; there are also two isolated populations in Northern Mexico.
It appears in the canyons of the Idaho Panhandle and more limitedly in the northeastern Olympic Mountains, which features some exceptionally large specimens, including one in diameter and tall. It can be found in the Cascade Range, mostly on the eastern slopes, from elevations of and liberally in the Rocky Mountains. It can also be found in the Monashee and Selkirk Mountains, as well as the highlands surrounding the Interior Plateau.
Ecology
Because transpiration is greatly reduced in small saplings while engulfed in snowpack, increased rates of transpiration in response to loss of snowpack, coupled with low sapwood water reserves and an extended period of soil frost in windswept areas, may prevent Engelmann spruce from regenerating in open areas both above and below the tree line. Both water uptake and water stored in roots appear to be critical for the survival of subalpine Engelmann spruce saplings that are exposed above the snowpack in later winter to early spring.
For exposed trees, the availability of soil water may be critical in late winter, when transpirational demands increase. Cuticular damage by windblown ice is probably more important at the tree line, but damage caused by desiccation is likely to be more important at lower elevations.
Despite wind damage, the species tends to grow taller than others at the tree line. It is shade tolerant, but not so much as subalpine fir, rendering it somewhat dependent on fires to outgrow competitors, although its thin bark and shallow roots make it vulnerable to fire. Spruce bark beetles attack the tree, being particularly deadly to groups which have stood for centuries. It is also susceptible to avalanches.
Although older spruce forests are not very useful to animals for forage, they can become so after fires, as they often burn completely, allowing many other plants, especially deciduous, to rise. Engelmann spruce-shaded streams are exploited by trout, and aphids produce galls which hang from the tree and look similar to cones when they dry out.
Subspecies and hybrids
Two geographical subspecies (treated as varieties by some authors, and as distinct species by others) occur:
Picea engelmannii subsp. engelmannii (Engelmann spruce). All of the range except as below.
Picea engelmannii subsp. mexicana (Mexican spruce). Two isolated populations on high mountains in northern Mexico, on the Sierra del Carmen in Coahuila (Sierra Madre Oriental) and on Cerro Mohinora in Chihuahua (Sierra Madre Occidental). Engelmann spruces of the Madrean sky islands mountains in the extreme southeast of Arizona and southwest of New Mexico also probably belong to this subspecies, though this is disputed.
The Engelmann spruce hybridises and intergrades extensively with the closely related white spruce (Picea glauca), found further north and east in the Rockies, and to a lesser extent with the closely related Sitka spruce where they meet on the western fringes of the Cascades.
Uses
Native Americans made various medicines from the resin and foliage.
Engelmann spruce is of economic importance for its wood, being light and fairly strong. It is harvested for paper-making and general construction. Wood from slow-grown trees at high elevation is especially prized for making soundboards for musical instruments such as acoustic guitars, harps, violins, and pianos.
Because it is odourless and has little resin, it has been used for food containers such as barrels. It is also used to a small extent as a Christmas tree.
Gallery
| Biology and health sciences | Pinaceae | Plants |
14411152 | https://en.wikipedia.org/wiki/Astrapotheria | Astrapotheria | Astrapotheria is an extinct order of South American and Antarctic hoofed mammals that existed from the late Paleocene to the Middle Miocene, . Astrapotheres were large, rhinoceros-like animals and have been called one of the most bizarre orders of mammals with an enigmatic evolutionary history.
The taxonomy of this order is not clear, but it may belong to Meridiungulata (along with Notoungulata, Litopterna, Pyrotheria and Xenungulata). In turn, Meridungulata is believed to belong to the extant superorder Laurasiatheria. Some scientists have regarded the astrapotheres (and sometimes the Meridiungulata as a whole) as members of the clade Atlantogenata. However, collagen and mitochondrial DNA sequence data analysed in 2015 places at least the notoungulates and litopterns firmly within Laurasiatheria, as a sister group to the perissodactyls.
Description
Their lophodont molars and tusk-like canines became extremely large and ever-growing in later astrapotheres. The upper molars lack an ectocingulum and are dominated by well-developed ectoloph and protoloph. Additional lophs formed in some derived taxa. They had lower molars with two cross-lophs, including a high protocristid, and eventually became almost selenodont. As a result, their dentition is similar to notoungulates, but it seems to have evolved independently. The cheek teeth are similar to rhinocerotoids, including similar microstructure, which indicate they had the same function.
Postcranially, astrapotheres are relatively robust and more or less graviportal but have slender long bones, most notably in the hindlegs, suggesting they were amphibious. In order to support their proboscises and large heads they had relatively long and massive necks in relation to the rest of the vertebral column. Their feet are pentadactyl with short and stout podial and metapodial bones. Most characteristic for the order are the flat astragalus, equipped with a short neck and a flat head, articulating with both the navicular and cuboid bones; and their calcaneus with its enlarged peroneal tubercle.
Three families are recognized: Eoastrapostylopidae from the late Paleocene, Trigonostylopidae from the Paleocene-Eocene, and Astrapotheriidae from the Eocene-Miocene. The Brazilian, Itaboraian Tetragonostylops and the Argentinian, Riochican Eoastrapostylops are the oldest astrapotheres. The latter, with its low-crowned and lophoselenodont cheek teeth, is considered the most primitive astrapothere. Trigonostylopids are distinct from other astrapotheres in their ear anatomy but are included in the order because of otherwise similar characters. Antarctodon is one of few eutherian mammals, as well as one of the last known terrestrial vertebrates, found in Antarctica.
The most famous member of the order is undoubtedly Astrapotherium, a long elephant-like animal that had lost its upper incisors and developed ever-growing canine tusks. They had lost their anterior premolars, resulting in a gap between their tusks and the hypsodont cheek teeth. The short and retracted nasal bones indicate a moderately developed tapir-like proboscis. The small Eocene Trigonostylops lacked such retracted nasals and probably also a proboscis. Other astrapotheriids, such as the Casamayoran Scaglia and Albertogaudrya, were between a sheep and a tapir in size and already the largest South American mammals.
Classification
There is no scientific consensus regarding the classification within Astrapotheria. For example, originally described Tetragonostylops as a trigonostylopid but Soria 1982 and 1984 transferred the genus to Astrapotheriidae and concluded that the remaining two genera in that family, Trigonostylops and Shecenia, form a basal collateral branch within Astrapotheriidae. According to , Trigonostylopidae (including Eoastrapostylopidae) is the stem group of Astrapotheriidae.
Astrapotheriidae
Albertogaudrya
Antarctodon
Astraponotus
Astrapothericulus
Astrapotherium
Comahuetherium
Granastrapotherium
Hilarcotherium
Liarthrus
Maddenia
Parastrapotherium
Scaglia
Uruguaytherium
Xenastrapotherium
Eoastrapostylopidae
Eoastrapostylops
Trigonostylopidae
Shecenia
Tetragonostylops
Trigonostylops
| Biology and health sciences | Mammals: General | Animals |
4086303 | https://en.wikipedia.org/wiki/Eczema%20herpeticum | Eczema herpeticum | Eczema herpeticum is a rare but severe and contagious disseminated infection that generally occurs at sites of skin damage produced by, for example, atopic dermatitis, burns, long-term usage of topical steroids or eczema. It is also known as Kaposi varicelliform eruption, Pustulosis varioliformis acute and Kaposi–Juliusberg dermatitis.
Some sources reserve the term "eczema herpeticum" when the cause is due to human herpes simplex virus, and the term "Kaposi varicelliform eruption" to describe the general presentation without specifying the virus.
This condition is most commonly caused by herpes simplex virus type 1 or 2, but may also be caused by coxsackievirus A16, or vaccinia virus. It appears as numerous umbilicated vesicles superimposed on healing atopic dermatitis. it is often accompanied by fever and lymphadenopathy. Eczema herpeticum can be life-threatening in babies.
Presentation
In addition to the skin, this infection affects multiple organs, including the eyes, brain, lung, and liver, and can be fatal.
Treatment
It can be treated with systemic antiviral drugs, such as aciclovir or valganciclovir. Foscarnet may also be used for immunocompromised host with Herpes simplex and acyclovir-resistant Herpes simplex.
Epidemiology
Even though the disease may develop at any age it is mostly present in childhood. Those who are affected typically have pre-existing cutaneous condition like atopic dermatitis.
History
Eczema herpeticum was first described by Hungarian dermatologist Moriz Kaposi in 1887. Fritz Juliusberg coined the term Pustulosis varioliformis acute in 1898. Eczema herpeticum is caused by Herpes simplex virus HSV1, the virus that causes cold sores; it can also be caused by other related viruses.
| Biology and health sciences | Viral diseases | Health |
4093158 | https://en.wikipedia.org/wiki/Strix%20%28bird%29 | Strix (bird) | Strix is a genus of owls in the typical owl family (Strigidae), one of the two generally accepted living families of owls, with the other being Tytonidae. Common names are earless owls or wood owls, though they are not the only owls without ear tufts, and "wood owl" is also used as a more generic name for forest-dwelling owls.
These are medium-sized to large, robustly built, powerful owls. They do not have ear tufts and most are highly nocturnal woodland birds. Most prey on small mammals, birds, and reptiles.
Most owls in the genus Strix can be distinguished from other genera of owls through their hooting vocalization and lack of visible ears.
The Latin genus name Strix referred to a mythical vampiric owl-monster believed to suck the blood of infants. Although the genus Strix was established for the earless owls by Linnaeus in 1758, many applied the term to other owls (namely the Tyto) until the late 19th century. This genus is closely related to the extinct Ornimegalonyx.
Taxonomy
The genus Strix was introduced by the Swedish naturalist Carl Linnaeus in 1758 in the tenth edition of his Systema Naturae. The type species is the tawny owl. The genus name is a Latin word meaning "owl".
Some Neotropical species were formerly classified in a separate genus, Ciccaba, which was eventually merged based on the placement of its type species, Strix huhula.
Species
The genus contains 22 species:
Spotted wood owl, S. seloputo
Mottled wood owl, S. ocellata
Brown wood owl, S. leptogrammica
Tawny owl, S. aluco
Maghreb owl, S. mauritanica
Himalayan owl, S. nivicolum
Desert owl, S. hadorami
Omani owl, S. butleri
Spotted owl, S. occidentalis
Barred owl, S. varia
Cinereous owl, S. sartorii
Fulvous owl, S. fulvescens
Rusty-barred owl, S. hylophila
Chaco owl, S. chacoensis
Rufous-legged owl, S. rufipes
Ural owl, S. uralensis
Great grey owl, S. nebulosa
African wood owl, S. woodfordii
Mottled owl, S. virgata
Black-and-white owl, S. nigrolineata
Black-banded owl, S. huhula
Rufous-banded owl, S. albitarsis
Fossil species
The genus Strix is well represented in the fossil record. Being a fairly generic type of strigid owl, they were probably the first truly modern Strigidae to evolve. However, whether several of the species usually placed in this genus indeed belong here is uncertain.
Generally accepted in Strix are:
S. dakota (Early Miocene of South Dakota, USA) – tentatively placed here
Strix sp. (Late Miocene of Nebraska, USA)
Strix sp. (Late Pliocene of Rębielice Królewski, Poland) apparently similar to the great grey owl
Strix intermedia (Early - Middle Pleistocene of EC Europe) – may be paleosubspecies of S. aluco
Strix brea (Late Pleistocene of SW North America) Now placed in its own genus. (See below)
Strix sp. (Late Pleistocene of Ladds, USA)
"Strix" wintershofensis (Early/Middle Miocene of Wintershof West, Germany) and "Strix" edwardsi (Middle Miocene of Grive-Saint-Alban, France), while being strigid owls, have not at present been reliably identified to genus; they might also belong into the European Ninox-like group.
"Strix" ignota (Middle Miocene of Sansan, France) is sometimes erroneously considered a nomen nudum, but this assumption is based on what appears to be a lapsus or misprint in a 1912 source. It may well belong into the present genus, but this requires confirmation.
"Strix" perpasta (Late Miocene – Early Pliocene of Gargano Peninsula, Italy) does not appear to belong into this genus either. It is sometimes considered a junior synonym of a brown fish-owl paleosubspecies.
UMMP V31030, a coracoid from Late Pliocene Rexroad Formation deposits of Kansas (USA), cannot be conclusively assigned to either the present genus or Bubo.
Extinct forms formerly in Strix:
"Strix" antiqua – now in Prosybris
"Strix" brea - now Oraristrix brea
"Strix" brevis – now in Intutula
"Strix" collongensis – now in Alasio
"Strix" melitensis and "Strix" sanctialbani – now in Tyto
"Strix" murivora – male of the Rodrigues scops owl
"Strix" newtoni and "Strix" sauzieri – male and female of the Mauritius scops owl
| Biology and health sciences | Strigiformes | Animals |
15923603 | https://en.wikipedia.org/wiki/Oil%20tanker | Oil tanker | An oil tanker, also known as a petroleum tanker, is a ship designed for the bulk transport of oil or its products. There are two basic types of oil tankers: crude tankers and product tankers. Crude tankers move large quantities of unrefined crude oil from its point of extraction to refineries. Product tankers, generally much smaller, are designed to move refined products from refineries to points near consuming markets.
Oil tankers are often classified by their size as well as their occupation. The size classes range from inland or coastal tankers of a few thousand metric tons of deadweight (DWT) to ultra-large crude carriers (ULCCs) of . Tankers move approximately of oil every year. Second only to pipelines in terms of efficiency, the average cost of transport of crude oil by tanker amounts to only US.
Some specialized types of oil tankers have evolved. One of these is the naval replenishment oiler, a tanker which can fuel a moving vessel. Combination ore-bulk-oil carriers and permanently moored floating storage units are two other variations on the standard oil tanker design. Oil tankers have been involved in a number of damaging and high-profile oil spills.
History
The technology of oil transportation has evolved alongside the oil industry. Although human use of oil reaches to prehistory, the first modern commercial exploitation dates back to James Young's manufacture of paraffin in 1850. In the early 1850s, oil began to be exported from Upper Burma, then a British colony. The oil was moved in earthenware vessels to the river bank where it was then poured into boat holds for transportation to Britain.
In the 1860s, Pennsylvania oil fields became a major supplier of oil, and a center of innovation after Edwin Drake had struck oil near Titusville, Pennsylvania. Break-bulk boats and barges were originally used to transport Pennsylvania oil in wooden barrels. But transport by barrel had several problems. The first problem was weight: they weighed , representing 20% of the total weight of a full barrel. Other problems with barrels were their expense, their tendency to leak, and the fact that they were generally used only once. The expense was significant: for example, in the early years of the Russian oil industry, barrels accounted for half the cost of petroleum production.
Early designs
In 1863, two sail-driven tankers were built on England's River Tyne. These were followed in 1873 by the first oil-tank steamer, Vaderland (Fatherland), which was built by Palmers Shipbuilding and Iron Company for Belgian owners. The vessel's use was curtailed by US and Belgian authorities citing safety concerns. By 1871, the Pennsylvania oil fields were making limited use of oil tank barges and cylindrical railroad tank-cars similar to those in use today.
Modern oil tankers
The modern oil tanker was developed in the period from 1877 to 1885. In 1876, Ludvig and Robert Nobel, brothers of Alfred Nobel, founded Branobel (short for Brothers Nobel) in Baku, Azerbaijan. It was, during the late 19th century, one of the largest oil companies in the world.
Ludvig was a pioneer in the development of early oil tankers. He first experimented with carrying oil in bulk on single-hulled barges. Turning his attention to self-propelled tankships, he faced a number of challenges. A primary concern was to keep the cargo and fumes well away from the engine room to avoid fires. Other challenges included allowing for the cargo to expand and contract due to temperature changes, and providing a method to ventilate the tanks.
The first successful oil tanker was Zoroaster, built by Sven Alexander Almqvist in Motala Verkstad, which carried its of kerosene cargo in two iron tanks joined by pipes. One tank was forward of the midships engine room and the other was aft. The ship also featured a set of 21 vertical watertight compartments for extra buoyancy. The ship had a length overall of , a beam of , and a draft of . Unlike later Nobel tankers, the Zoroaster design was built small enough to sail from Sweden to the Caspian by way of the Baltic Sea, Lake Ladoga, Lake Onega, the Rybinsk and Mariinsk Canals and the Volga River. The aft and the stern was put together and then dismantled to make room for the mid-section as the Caspian Sea was reached.
In 1883, oil tanker design took a large step forward. Working for the Nobel company, British engineer Colonel Henry F. Swan designed a set of three Nobel tankers. Instead of one or two large holds, Swan's design used several holds which spanned the width, or beam, of the ship. These holds were further subdivided into port and starboard sections by a longitudinal bulkhead. Earlier designs suffered from stability problems caused by the free surface effect, where oil sloshing from side to side could cause a ship to capsize. But this approach of dividing the ship's storage space into smaller tanks virtually eliminated free-surface problems. This approach, almost universal today, was first used by Swan in the Nobel tankers Blesk, Lumen, and Lux.
Others point to , another design of Colonel Swan, as being the first modern oil tanker. It adopted the best practices from previous oil tanker designs to create the prototype for all subsequent vessels of the type. It was the first dedicated steam-driven ocean-going tanker in the world and was the first ship in which oil could be pumped directly into the vessel hull instead of being loaded in barrels or drums. It was also the first tanker with a horizontal bulkhead; its features included cargo valves operable from the deck, cargo main piping, a vapor line, cofferdams for added safety, and the ability to fill a ballast tank with seawater when empty of cargo. The ship was built in Britain, and was purchased by Wilhelm Anton Riedemann, an agent for the Standard Oil Company along with several of her sister ships. After Glückauf was lost in 1893 after being grounded in fog, Standard Oil purchased the sister ships.
Asian trade
The 1880s also saw the beginnings of the Asian oil trade. The idea that led to moving Russian oil to the Far East via the Suez Canal was the brainchild of two men: importer Marcus Samuel and shipowner/broker Fred Lane. Prior bids to move oil through the canal had been rejected by the Suez Canal Company as being too risky. Samuel approached the problem a different way: asking the company for the specifications of a tanker it would allow through the canal.
Armed with the canal company's specifications, Samuel ordered three tankers from William Gray & Company in northern England. Named , Conch and Clam, each had a capacity of 5,010 long tons of deadweight. These three ships were the first tankers of the Tank Syndicate, forerunner of today's Royal Dutch Shell company.
With facilities prepared in Jakarta, Singapore, Bangkok, Saigon, Hong Kong, Shanghai, and Kobe, the fledgling Shell company was ready to become Standard Oil's first challenger in the Asian market. On August 24, 1892, Murex became the first tanker to pass through the Suez Canal. By the time Shell merged with Royal Dutch Petroleum in 1907, the company had 34 steam-driven oil tankers, compared to Standard Oil's four case-oil steamers and 16 sailing tankers.
The supertanker era
Until 1956, tankers were designed to be able to navigate the Suez Canal. This size restriction became much less of a priority after the closing of the canal during the Suez Crisis of 1956. Forced to move oil around the Cape of Good Hope, shipowners realized that bigger tankers were the key to more efficient transport. While a typical T2 tanker of the World War II era was long and had a capacity of , the ultra-large crude carriers (ULCC) built in the 1970s were over long and had a capacity of . Several factors encouraged this growth. Hostilities in the Middle East which interrupted traffic through the Suez Canal contributed, as did nationalization of Middle East oil refineries. Fierce competition among shipowners also played a part. But apart from these considerations is a simple economic advantage: the larger an oil tanker is, the more cheaply it can move crude oil, and the better it can help meet growing demands for oil.
In 1955 the world's largest supertanker was and : SS Spyros Niarchos launched that year by Vickers Armstrongs Shipbuilders Ltd in England for Greek shipping magnate Stavros Niarchos.
In 1958 United States shipping magnate Daniel K. Ludwig broke the record of 100,000 long tons of heavy displacement. His Universe Apollo displaced 104,500 long tons, a 23% increase from the previous record-holder, Universe Leader which also belonged to Ludwig. The first tanker over 100,000 dwt built in Europe was the British Admiral. The ship was launched at Barrow-in-Furness in 1965 by Elizabeth II.
The world's largest supertanker was built in 1979 at the Oppama shipyard by Sumitomo Heavy Industries, Ltd., named Seawise Giant. This ship was built with a capacity of , a length overall of and a draft of . She had 46 tanks, of deck, and at her full load draft, could not navigate the English Channel.
Seawise Giant was renamed Happy Giant in 1989, Jahre Viking in 1991, and Knock Nevis in 2004 (when she was converted into a permanently moored storage tanker). In 2009 she was sold for the last time, renamed Mont, and scrapped.
As of 2011, the world's two largest working supertankers are the s TI Europe and TI Oceania. These ships were built in 2002 and 2003 as Hellespont Alhambra and Hellespont Tara for the Greek Hellespont Steamship Corporation. Hellespont sold these ships to Overseas Shipholding Group and Euronav in 2004. Each of the sister ships has a capacity of over , a length overall of and a cargo capacity of . They were the first ULCCs to be double-hulled. To differentiate them from smaller ULCCs, these ships are sometimes given the V-Plus size designation.
With the exception of the pipeline, the tanker is the most cost-effective way to move oil today. Worldwide, tankers carry some annually, and the cost of transportation by tanker amounts to only US$0.02 per gallon at the pump.
Size categories
In 1954, Shell Oil developed the "average freight rate assessment" (AFRA) system which classifies tankers of different sizes. To make it an independent instrument, Shell consulted the London Tanker Brokers' Panel (LTBP). At first, they divided the groups as General Purpose for tankers under ; Medium Range for ships between 25,000 and 45,000 DWT and Long Range for the then-enormous ships that were larger than 45,000 DWT. The ships became larger during the 1970s, which prompted rescaling.
The system was developed for tax reasons as the tax authorities wanted evidence that the internal billing records were correct. Before the New York Mercantile Exchange started trading crude oil futures in 1983, it was difficult to determine the exact price of oil, which could change with every contract. Shell and BP, the first companies to use the system, abandoned the AFRA system in 1983, later followed by the US oil companies. However, the system is still used today. Besides that, there is the flexible market scale, which takes typical routes and lots of .
Merchant oil tankers carry a wide range of hydrocarbon liquids ranging from crude oil to refined petroleum products. Crude carriers are among the largest, ranging from 55,000 DWT Panamax-sized vessels to ultra-large crude carriers (ULCCs) of over 440,000 DWT.
Smaller tankers, ranging from well under 10,000 DWT to 80,000 DWT Panamax vessels, generally carry refined petroleum products, and are known as product tankers. The smallest tankers, with capacities under 10,000 DWT generally work near-coastal and inland waterways. Although they were in the past, ships of the smaller Aframax and Suezmax classes are no longer regarded as supertankers.
VLCC and ULCC
"Supertankers" are the largest oil tankers, and the largest mobile man-made structures. They include very large and ultra-large crude carriers (VLCCs and ULCCs – see above) with capacities over 250,000 DWT. These ships can transport of oil/318,000 metric tons. By way of comparison, the United Kingdom consumed about of oil per day in 2009. ULCCs commissioned in the 1970s were the largest vessels ever built, but have all now been scrapped. A few newer ULCCs remain in service, none of which are more than 400 meters long.
Because of their size, supertankers often cannot enter port fully loaded. These ships can take on their cargo at offshore platforms and single-point moorings. On the other end of the journey, they often pump their cargo off to smaller tankers at designated lightering points off-coast. Supertanker routes are typically long, requiring them to stay at sea for extended periods, often around seventy days at a time.
Chartering
The act of hiring a ship to carry cargo is called chartering. (The contract itself is known as a charter party.) Tankers are hired by four types of charter agreements: the voyage charter, the time charter, the bareboat charter, and contract of affreightment. In a voyage charter the charterer rents the vessel from the loading port to the discharge port. In a time charter the vessel is hired for a set period of time, to perform voyages as the charterer directs. In a bareboat charter the charterer acts as the ship's operator and manager, taking on responsibilities such as providing the crew and maintaining the vessel. Finally, in a contract of affreightment or COA, the charterer specifies a total volume of cargo to be carried in a specific time period and in specific sizes, for example a COA could be specified as of JP-5 in a year's time in shipments.
One of the key aspects of any charter party is the freight rate, or the price specified for carriage of cargo. The freight rate of a tanker charter party is specified in one of four ways: by a lump sum rate, by rate per ton, by a time charter equivalent rate, or by Worldscale rate. In a lump sum rate arrangement, a fixed price is negotiated for the delivery of a specified cargo, and the ship's owner/operator is responsible to pay for all port costs and other voyage expenses. Rate per ton arrangements are used mostly in chemical tanker chartering, and differ from lump sum rates in that port costs and voyage expenses are generally paid by the charterer. Time charter arrangements specify a daily rate, and port costs and voyage expenses are also generally paid by the charterer.
The Worldwide Tanker Normal Freight Scale, often referred to as Worldscale, is established and governed jointly by the Worldscale Associations of London and New York. Worldscale establishes a baseline price for carrying a metric ton of product between any two ports in the world. In Worldscale negotiations, operators and charterers will determine a price based on a percentage of the Worldscale rate. The baseline rate is expressed as WS 100. If a given charter party settled on 85% of the Worldscale rate, it would be expressed as WS 85. Similarly, a charter party set at 125% of the Worldscale rate would be expressed as WS 125.
Recent markets
The market is affected by a wide variety of variables such as the supply and demand of oil as well as the supply and demand of oil tankers. Some particular variables include winter temperatures, excess tanker tonnage, supply fluctuations in the Persian Gulf, and interruptions in refinery services.
In 2006, time-charters tended towards long term. Of the time charters executed in that year, 58% were for a period of 24 or more months, 14% were for periods of 12 to 24 months, 4% were from 6 to 12 months, and 24% were for periods of less than 6 months.
From 2003, the demand for new ships started to grow, resulting in 2007 in a record breaking order backlog for shipyards, exceeding their capacity with rising newbuilding prices as a result. This resulted in a glut of ships when demand dropped due to a weakened global economy and dramatically reduced demand in the United States. The charter rate for very large crude carriers, which carry two million barrels of oil, had peaked at $309,601 per day in 2007 but had dropped to $7,085 per day by 2012, far below the operating costs of these ships. As a result, several tanker operators laid up their ships. Prices rose significantly in 2015 and early 2016, but delivery of new tankers was projected to keep prices in check.
Owners of large oil tanker fleets include Teekay Corporation, A P Moller Maersk, DS Torm, Frontline, MOL Tankship Management, Overseas Shipholding Group, and Euronav.
Fleet characteristics
In 2005, oil tankers made up 36.9% of the world's fleet in terms of deadweight tonnage. The world's total oil tankers deadweight tonnage has increased from in 1970 to in 2005. The combined deadweight tonnage of oil tankers and bulk carriers, represents 72.9% of the world's fleet.
Cargo movement
In 2005, 2.42 billion metric tons of oil were shipped by tanker. 76.7% of this was crude oil, and the rest consisted of refined petroleum products. This amounted to 34.1% of all seaborne trade for the year. Combining the amount carried with the distance it was carried, oil tankers moved 11,705 billion metric-ton-miles of oil in 2005.
By comparison, in 1970 1.44 billion metric tons of oil were shipped by tanker. This amounted to 34.1% of all seaborne trade for that year. In terms of amount carried and distance carried, oil tankers moved 6,487 billion metric-ton-miles of oil in 1970.
The United Nations also keeps statistics about oil tanker productivity, stated in terms of metric tons carried per metric ton of deadweight as well as metric-ton-miles of carriage per metric ton of deadweight. In 2005, for each of oil tankers, 6.7 metric tons of cargo was carried. Similarly, each of oil tankers was responsible for 32,400 metric-ton miles of carriage.
The main loading ports in 2005 were located in Western Asia, Western Africa, North Africa, and the Caribbean, with 196.3, 196.3, 130.2 and 246.6 million metric tons of cargo loaded in these regions. The main discharge ports were located in North America, Europe, and Japan with 537.7, 438.4, and 215.0 million metric tons of cargo discharged in these regions.
Flag states
International law requires that every merchant ship be registered in a country, called its flag state. A ship's flag state exercises regulatory control over the vessel and is required to inspect it regularly, certify the ship's equipment and crew, and issue safety and pollution prevention documents. As of 2007, the United States Central Intelligence Agency statistics count 4,295 oil tankers of or greater worldwide. Panama was the world's largest flag state for oil tankers, with 528 of the vessels in its registry. Six other flag states had more than 200 registered oil tankers: Liberia (464), Singapore (355), China (252), Russia (250), the Marshall Islands (234) and the Bahamas (209). The Panamanian, Liberian, Marshallese and Bahamian flags are open registries and considered by the International Transport Workers' Federation to be flags of convenience. By comparison, the United States and the United Kingdom only had 59 and 27 registered oil tankers, respectively.
Vessel life cycle
In 2005, the average age of oil tankers worldwide was 10 years. Of these, 31.6% were under 4 years old and 14.3% were over 20 years old. In 2005, 475 new oil tankers were built, accounting for . The average size for these new tankers was . Nineteen of these were VLCC size, 19 were Suezmax, 51 were Aframax, and the rest were smaller designs. By comparison, , , and worth of oil tanker capacity was built in 1980, 1990, and 2000 respectively.
Ships are generally removed from the fleet through a process known as scrapping. Ship-owners and buyers negotiate scrap prices based on factors such as the ship's empty weight (called light ton displacement or LDT) and prices in the scrap metal market. In 1998, almost 700 ships went through the scrapping process at shipbreakers in places such as Gadani, Alang and Chittagong. In 2004 and 2005, and respectively of oil tankers were scrapped. Between 2000 and 2005, the capacity of oil tankers scrapped each year has ranged between and . In this same timeframe, tankers have accounted for between 56.5% and 90.5% of the world's total scrapped ship tonnage. In this period the average age of scrapped oil tankers has ranged from 26.9 to 31.5 years.
Vessel pricing
In 2005, the price for new oil tankers in the , , and ranges were $43 million, $58 million, and $120 million respectively. In 1985 these vessels would have cost $18 million, $22 million, and $47 million respectively.
Oil tankers are often sold second hand. In 2005, worth of oil tankers were sold used. Some representative prices for that year include $42.5 million for a tanker, $60.7 million for a , $73 million for a , and $116 million for tanker. For a concrete example, in 2006, Bonheur subsidiary First Olsen paid $76.5 million for Knock Sheen, a 159,899 DWT tanker.
The cost of operating the largest tankers, the Very Large Crude Carriers, is currently between $10,000 and $12,000 per day.
Current structural design
Oil tankers generally have from 8 to 12 tanks. Each tank is split into two or three independent compartments by fore-and-aft bulkheads. The tanks are numbered with tank one being the forwardmost. Individual compartments are referred to by the tank number and the athwartships position, such as "one port", "three starboard", or "six center".
A cofferdam is a small space left open between two bulkheads, to give protection from heat, fire, or collision. Tankers generally have cofferdams forward and aft of the cargo tanks, and sometimes between individual tanks. A pumproom houses all the pumps connected to a tanker's cargo lines. Some larger tankers have two pumprooms. A pumproom generally spans the total breadth of the ship.
Hull designs
A major component of tanker architecture is the design of the hull or outer structure. A tanker with a single outer shell between the product and the ocean is said to be "single-hulled". Most newer tankers are "double hulled", with an extra space between the hull and the storage tanks. Hybrid designs such as "double-bottom" and "double-sided" combine aspects of single and double-hull designs. All single-hulled tankers around the world will be phased out by 2026, in accordance with the International Convention for the Prevention of Pollution from Ships, 1973 (MARPOL). The United Nations has decided to phase out single hull oil tankers by 2010.
In 1998, the Marine Board of the National Academy of Sciences conducted a survey of industry experts regarding the pros and cons of double-hull design. Some of the advantages of the double-hull design that were mentioned include ease of ballasting in emergency situations, reduced practice of saltwater ballasting in cargo tanks decreases corrosion, increased environmental protection, cargo discharge is quicker, more complete and easier, tank washing is more efficient, and better protection in low-impact collisions and grounding.
The same report lists the following as some drawbacks to the double-hull design, including higher build costs, greater operating expenses (e.g. higher canal and port tariffs), difficulties in ballast tank ventilation, the fact that ballast tanks need continuous monitoring and maintenance, increased transverse free surface, the greater number of surfaces to maintain, the risk of explosions in double-hull spaces if a vapor detection system not fitted, and that cleaning ballast tanks is more difficult for double hull ships.
In all, double-hull tankers are said to be safer than a single-hull in a grounding incident, especially when the shore is not very rocky. The safety benefits are less clear on larger vessels and in cases of high speed impact.
Although double-hull design is superior in low energy casualties and prevents spillage in small casualties, in high energy casualties where both hulls are breached, oil can spill through the double-hull and into the sea and spills from a double-hull tanker can be significantly higher than designs like the mid-deck tanker, the Coulombi Egg Tanker and even a pre-MARPOL tanker, as the last one has a lower oil column and reaches hydrostatic balance sooner.
Inert gas system
An oil tanker's inert gas system is one of the most important parts of its design. Fuel oil itself is very difficult to ignite, but its hydrocarbon vapors are explosive when mixed with air in certain concentrations. The purpose of the system is to create an atmosphere inside tanks in which the hydrocarbon oil vapors cannot burn.
As inert gas is introduced into a mixture of hydrocarbon vapors and air, it increases the lower flammable limit or lowest concentration at which the vapors can be ignited. At the same time it decreases the upper flammable limit or highest concentration at which the vapors can be ignited. When the total concentration of oxygen in the tank decreases to about 11%, the upper and lower flammable limits converge and the flammable range disappears.
Inert gas systems deliver air with an oxygen concentration of less than 5% by volume. As a tank is pumped out, it is filled with inert gas and kept in this safe state until the next cargo is loaded. The exception is in cases when the tank must be entered. Safely gas-freeing a tank is accomplished by purging hydrocarbon vapors with inert gas until the hydrocarbon concentration inside the tank is under about 1%. Thus, as air replaces the inert gas, the concentration cannot rise to the lower flammable limit and is safe.
Cargo operations
Operations aboard oil tankers are governed by an established body of best practices and a large body of international law. Cargo can be moved on or off of an oil tanker in several ways. One method is for the ship to moor alongside a pier, connect with cargo hoses or marine loading arms. Another method involves mooring to offshore buoys, such as a single point mooring, and making a cargo connection via underwater cargo hoses. A third method is by ship-to-ship transfer, also known as lightering. In this method, two ships come alongside in open sea and oil is transferred manifold to manifold via flexible hoses. Lightering is sometimes used where a loaded tanker is too large to enter a specific port.
Pre-transfer preparation
Prior to any transfer of cargo, the chief officer must develop a transfer plan detailing specifics of the operation such as how much cargo will be moved, which tanks will be cleaned, and how the ship's ballasting will change. The next step before a transfer is the pretransfer conference. The pretransfer conference covers issues such as what products will be moved, the order of movement, names and titles of key people, particulars of shipboard and shore equipment, critical states of the transfer, regulations in effect, emergency and spill-containment procedures, watch and shift arrangements, and shutdown procedures.
After the conference is complete, the person in charge on the ship and the person in charge of the shore installation go over a final inspection checklist. In the United States, the checklist is called a Declaration of Inspection or DOI. Outside the US, the document is called the "Ship/Shore Safety Checklist." Items on the checklist include proper signals and signs are displayed, secure mooring of the vessel, choice of language for communication, securing of all connections, that emergency equipment is in place, and that no repair work is taking place.
Loading cargo
Loading an oil tanker consists primarily of pumping cargo into the ship's tanks. As oil enters the tank, the vapors inside the tank must be somehow expelled. Depending on local regulations, the vapors can be expelled into the atmosphere or discharged back to the pumping station by way of a vapor recovery line. It is also common for the ship to move water ballast during the loading of cargo to maintain proper trim.
Loading starts slowly at a low pressure to ensure that equipment is working correctly and that connections are secure. Then a steady pressure is achieved and held until the "topping-off" phase when the tanks are nearly full. Topping off is a very dangerous time in handling oil, and the procedure is handled particularly carefully. Tank-gauging equipment is used to tell the person in charge how much space is left in the tank, and all tankers have at least two independent methods for tank-gauging. As the tanker becomes full, crew members open and close valves to direct the flow of product and maintain close communication with the pumping facility to decrease and finally stop the flow of liquid.
Unloading cargo
The process of moving oil off of a tanker is similar to loading, but has some key differences. The first step in the operation is following the same pretransfer procedures as used in loading. When the transfer begins, it is the ship's cargo pumps that are used to move the product ashore. As in loading, the transfer starts at low pressure to ensure that equipment is working correctly and that connections are secure. Then a steady pressure is achieved and held during the operation. While pumping, tank levels are carefully watched and key locations, such as the connection at the cargo manifold and the ship's pumproom are constantly monitored. Under the direction of the person in charge, crew members open and close valves to direct the flow of product and maintain close communication with the receiving facility to decrease and finally stop the flow of liquid.
Tank cleaning
Tanks must be cleaned from time to time for various reasons. One reason is to change the type of product carried inside a tank. Also, when tanks are to be inspected or maintenance must be performed within a tank, it must be not only cleaned, but made gas-free.
On most crude-oil tankers, a special crude oil washing (COW) system is part of the cleaning process. The COW system circulates part of the cargo through the fixed tank-cleaning system to remove wax and asphaltic deposits. Tanks that carry less viscous cargoes are washed with water. Fixed and portable automated tank cleaning machines, which clean tanks with high-pressure water jets, are widely used. Some systems use rotating high-pressure water jets to spray hot water on all the internal surfaces of the tank. As the spraying takes place, the liquid is pumped out of the tank.
After a tank is cleaned, provided that it is going to be prepared for entry, it will be purged. Purging is accomplished by pumping inert gas into the tank until hydrocarbons have been sufficiently expelled. Next the tank is gas freed which is usually accomplished by blowing fresh air into the space with portable air powered or water powered air blowers. "Gas freeing" brings the oxygen content of the tank up to 20.8% O2. The inert gas buffer between fuel and oxygen atmospheres ensures they are never capable of ignition. Specially trained personnel monitor the tank's atmosphere, often using hand-held gas indicators which measure the percentage of hydrocarbons present. After a tank is gas-free, it may be further hand-cleaned in a manual process known as mucking. Mucking requires protocols for entry into confined spaces, protective clothing, designated safety observers, and possibly the use of airline respirators.
Special-use oil tankers
Some sub-types of oil tankers have evolved to meet specific military and economic needs. These sub-types include naval replenishment ships, oil-bulk-ore combination carriers, floating storage and offloading units (FSOs) and floating production storage and offloading units (FPSOs).
Replenishment ships
Replenishment ships, known as oilers in the United States and fleet tankers in Commonwealth countries, are ships that can provide oil products to naval vessels while on the move. This process, called underway replenishment, extends the length of time a naval vessel can stay at sea, as well as her effective range. Prior to underway replenishment, naval vessels had to enter a port or anchor to take on fuel. In addition to fuel, replenishment ships may also deliver water, ammunition, rations, stores and personnel.
Ore-bulk-oil carriers
An ore-bulk-oil carrier, also known as combination carrier or OBO, is a ship designed to be capable of carrying wet or dry bulk cargoes. This design was intended to provide flexibility in two ways. Firstly, an OBO would be able to switch between the dry and wet bulk trades based on market conditions. Secondly, an OBO could carry oil on one leg of a voyage and return carrying dry bulk, reducing the number of unprofitable ballast voyages it would have to make.
In practice, the flexibility which the OBO design allows has gone largely unused, as these ships tend to specialize in either the liquid or dry bulk trade. Also, these ships have endemic maintenance problems. On one hand, due to a less specialized design, an OBO suffers more from wear and tear during dry cargo onload than a bulker. On the other hand, components of the liquid cargo system, from pumps to valves to piping, tend to develop problems when subjected to periods of disuse. These factors have contributed to a steady reduction in the number of OBO ships worldwide since the 1970s.
One of the more famous OBOs was of which in September 1980 became the largest British ship ever lost at sea. It sank in a Pacific typhoon while carrying a cargo of iron ore from Canada to Japan.
Floating storage units
Floating storage and offloading units (FSO) are used worldwide by the offshore oil industry to receive oil from nearby platforms and store it until it can be offloaded onto oil tankers. A similar system, the floating production storage and offloading unit (FPSO), has the ability to process the product while it is on board. These floating units reduce oil production costs and offer mobility, large storage capacity, and production versatility.
FPSO and FSOs are often created out of old, stripped-down oil tankers, but can be made from new-built hulls; Shell España first used a tanker as an FPSO in August 1977. An example of an FSO that used to be an oil tanker is the Knock Nevis. These units are usually moored to the seabed through a spread mooring system. A turret-style mooring system can be used in areas prone to severe weather. This turret system lets the unit rotate to minimize the effects of sea-swell and wind.
Pollution
Oil spills have devastating effects on the environment. Crude oil contains polycyclic aromatic hydrocarbons (PAHs) which are very difficult to clean up, and last for years in the sediment and marine environment. Marine species constantly exposed to PAHs can exhibit developmental problems, susceptibility to disease, and abnormal reproductive cycles.
By the sheer amount of oil carried, modern oil tankers can be a threat to the environment. As discussed above, a VLCC tanker can carry of crude oil. This is about eight times the amount spilled in the widely known Exxon Valdez incident. In this spill, the ship ran aground and dumped of oil into the ocean in March 1989. Despite efforts of scientists, managers, and volunteers, over 400,000 seabirds, about 1,000 sea otters, and immense numbers of fish were killed. Considering the volume of oil carried by sea, however, tanker owners' organizations often argue that the industry's safety record is excellent, with only a tiny fraction of a percentage of oil cargoes carried ever being spilled. The International Association of Independent Tanker Owners has observed that "accidental oil spills this decade have been at record low levels—one third of the previous decade and one tenth of the 1970s—at a time when oil transported has more than doubled since the mid 1980s."
Oil tankers are only one source of oil spills. According to the United States Coast Guard, 35.7% of the volume of oil spilled in the United States from 1991 to 2004 came from tank vessels (ships/barges), 27.6% from facilities and other non-vessels, 19.9% from non-tank vessels, 9.3% from pipelines, and 7.4% from mystery spills. Only 5% of the actual spills came from oil tankers, while 51.8% came from other kinds of vessels. The detailed statistics for 2004 show tank vessels responsible for somewhat less than 5% of the number of total spills but more than 60% of the volume. Tanker spills are much more rare and much more serious than spills from non-tank vessels.
The International Tanker Owners Pollution Federation has tracked 9,351 accidental spills that have occurred since 1974. According to this study, most spills result from routine operations such as loading cargo, discharging cargo, and taking on fuel oil. 91% of the operational oil spills are small, resulting in less than 7 metric tons per spill. On the other hand, spills resulting from accidents like collisions, groundings, hull failures, and explosions are much larger, with 84% of these involving losses of over 700 metric tons.
Following the Exxon Valdez spill, the United States passed the Oil Pollution Act of 1990 (OPA-90), which excluded single-hull tank vessels of 5,000 gross tons or more from US waters from 2010 onward, apart from those with a double bottom or double sides, which may be permitted to trade to the United States through 2015, depending on their age. Following the sinkings of (1999) and Prestige (2002), the European Union passed its own stringent anti-pollution packages (known as Erika I, II, and III), which also require all tankers entering its waters to be double-hulled by 2010. The Erika packages are controversial because they introduced the new legal concept of "serious negligence".
Air pollution
Large ships are often run on low quality fuel oils, such as bunker oil, which is highly polluting and has been shown to be a health risk.
| Technology | Naval transport | null |
15925628 | https://en.wikipedia.org/wiki/Hepatitis%20B | Hepatitis B | Hepatitis B is an infectious disease caused by the hepatitis B virus (HBV) that affects the liver; it is a type of viral hepatitis. It can cause both acute and chronic infection.
Many people have no symptoms during an initial infection. For others, symptoms may appear 30 to 180 days after becoming infected and can include a rapid onset of sickness with nausea, vomiting, yellowish skin, fatigue, yellow urine, and abdominal pain. Symptoms during acute infection typically last for a few weeks, though some people may feel sick for up to six months. Deaths resulting from acute stage HBV infections are rare. An HBV infection lasting longer than six months is usually considered chronic. The likelihood of developing chronic hepatitis B is higher for those who are infected with HBV at a younger age. About 90% of those infected during or shortly after birth develop chronic hepatitis B, while less than 10% of those infected after the age of five develop chronic cases. Most of those with chronic disease have no symptoms; however, cirrhosis and liver cancer eventually develop in about 25% of those with chronic HBV.
The virus is transmitted by exposure to infectious blood or body fluids. In areas where the disease is common, infection around the time of birth or from contact with other people's blood during childhood are the most frequent methods by which hepatitis B is acquired. In areas where the disease is rare, intravenous drug use and sexual intercourse are the most frequent routes of infection. Other risk factors include working in healthcare, blood transfusions, dialysis, living with an infected person, travel in countries with high infection rates, and living in an institution. Tattooing and acupuncture led to a significant number of cases in the 1980s; however, this has become less common with improved sterilization. The viruses cannot be spread by holding hands, sharing eating utensils, kissing, hugging, coughing, sneezing, or breastfeeding. The infection can be diagnosed 30 to 60 days after exposure. The diagnosis is usually confirmed by testing the blood for parts of the virus and for antibodies against the virus. It is one of five main hepatitis viruses: A, B, C, D, and E. During an initial infection, care is based on a person's symptoms. In those who develop chronic disease, antiviral medication such as tenofovir or interferon may be useful; however, these drugs are expensive. Liver transplantation is sometimes recommended for cases of cirrhosis or hepatocellular carcinoma.
Hepatitis B infection has been preventable by vaccination since 1982. As of 2022, the hepatitis B vaccine is between 98% and 100% effective in preventing infection. The vaccine is administered in several doses; after an initial dose, two or three more vaccine doses are required at a later time for full effect. The World Health Organization (WHO) recommends infants receive the vaccine within 24 hours after birth when possible. National programs have made the hepatitis B vaccine available for infants in 190 countries as of the end of 2021. To further prevent infection, the WHO recommends testing all donated blood for hepatitis B before using it for transfusion. Using antiviral prophylaxis to prevent mother-to-child transmission is also recommended, as is following safe sex practices, including the use of condoms. In 2016, the WHO set a goal of eliminating viral hepatitis as a threat to global public health by 2030. Achieving this goal would require the development of therapeutic treatments to cure chronic hepatitis B, as well as preventing its transmission and using vaccines to prevent new infections.
An estimated 296 million people, or 3.8% of the global population, had chronic hepatitis B infections as of 2019. Another 1.5 million developed acute infections that year, and 820,000 deaths occurred as a result of HBV. Cirrhosis and liver cancer are responsible for most HBV-related deaths. The disease is most prevalent in Africa (affecting 7.5% of the continent's population) and in the Western Pacific region (5.9%). Infection rates are 1.5% in Europe and 0.5% in the Americas. According to some estimates, about a third of the world's population has been infected with hepatitis B at one point in their lives. Hepatitis B was originally known as "serum hepatitis".
Signs and symptoms
Acute infection with virus is associated with acute viral hepatitis, an illness that begins with general ill-health, loss of appetite, nausea, vomiting, body aches, mild fever, and dark urine, and then progresses to development of jaundice. The illness lasts for a few weeks and then gradually improves in most affected people. A few people may have a more severe form of liver disease known as fulminant hepatic failure and may die as a result. The infection may be entirely asymptomatic and may go unrecognized.
Chronic infection with virus may be asymptomatic or may be associated with chronic inflammation of the liver (chronic hepatitis), leading to cirrhosis over a period of several years. This type of infection dramatically increases the incidence of hepatocellular carcinoma (HCC; liver cancer). Across Europe, hepatitis B and C cause approximately 50% of hepatocellular carcinomas. Chronic carriers are encouraged to avoid consuming alcohol as it increases their risk for cirrhosis and liver cancer. virus has been linked to the development of membranous glomerulonephritis (MGN).
Symptoms outside of the liver are present in 1–10% of HBV-infected people and include serum-sickness–like syndrome, acute necrotizing vasculitis (polyarteritis nodosa), membranous glomerulonephritis, and papular acrodermatitis of childhood (Gianotti–Crosti syndrome). The serum-sickness–like syndrome occurs in the setting of acute , often preceding the onset of jaundice. The clinical features are fever, skin rash, and polyarteritis. The symptoms often subside shortly after the onset of jaundice but can persist throughout the duration of acute . About 30–50% of people with acute necrotizing vasculitis (polyarteritis nodosa) are HBV carriers. HBV-associated nephropathy has been described in adults but is more common in children. Membranous glomerulonephritis is the most common form. Other immune-mediated hematological disorders, such as essential mixed cryoglobulinemia and aplastic anemia have been described as part of the extrahepatic manifestations of HBV infection, but their association is not as well-defined; therefore, they probably should not be considered etiologically linked to HBV.
Cause
Transmission
Transmission of virus results from exposure to infectious blood or body fluids containing blood. HBV is 50 to 100 times more infectious than human immunodeficiency virus (HIV). HBV can be transmitted through several routes of infection. In vertical transmission, HBV is passed from mother to child (MTCT) during childbirth. Without intervention, a mother who is positive for HBsAg has a 20% risk of passing the infection to her offspring at the time of birth. This risk is as high as 90% if the mother is also positive for HBeAg.
Early life horizontal transmission can occur through bites, lesions, certain sanitary habits, or other contact with secretions or saliva containing HBV. Adult horizontal transmission is known to occur through sexual contact, blood transfusions and transfusion with other human blood products, re-use of contaminated needles and syringes. Breastfeeding after proper immunoprophylaxis does not appear to contribute to mother-to-child-transmission (MTCT) of HBV.
Virology
Structure
virus (HBV) is a member of the hepadnavirus family. The virus particle (virion) consists of an outer lipid envelope and an icosahedral nucleocapsid core composed of core protein. These virions are 30–42 nm in diameter. The nucleocapsid encloses the viral DNA and a DNA polymerase that has reverse transcriptase activity. The outer envelope contains embedded proteins that are involved in viral binding of, and entry into, susceptible cells. The virus is one of the smallest enveloped animal viruses. The 42 nm virions, which are capable of infecting liver cells known as hepatocytes, are referred to as "Dane particles". In addition to the Dane particles, filamentous and spherical bodies lacking a core can be found in the serum of infected individuals. These particles are not infectious and are composed of the lipid and protein that forms part of the surface of the virion, which is called the surface antigens (HBsAg), and is produced in excess during the life cycle of the virus.
Genome
The genome of HBV is made of circular DNA, but it is unusual because the DNA is not fully double-stranded. One end of the full length strand is linked to the HBV DNA polymerase. The genome is 3020–3320 nucleotides long (for the full-length strand) and 1700–2800 nucleotides long (for the short length-strand). The negative-sense (non-coding) is complementary to the viral mRNA. The viral DNA is found in the nucleus soon after infection of the cell. The partially double-stranded DNA is rendered fully double-stranded by completion of the (+) sense strand and removal of a protein molecule from the (−) sense strand and a short sequence of RNA from the (+) sense strand. Non-coding bases are removed from the ends of the (−) sense strand and the ends are rejoined. There are four known genes encoded by the genome, called C, X, P, and S. The core protein is coded for by gene C (HBcAg), and its start codon is preceded by an upstream in-frame AUG start codon from which the pre-core protein is produced. HBeAg is produced by proteolytic processing of the pre-core protein. In some rare strains of the virus known as hepatitis B virus precore mutants, no HBeAg is present.
The DNA polymerase is encoded by gene P. Gene S is the gene that codes for the surface antigen (HBsAg). The HBsAg gene is one long open reading frame but contains three in frame "start" (ATG) codons that divide the gene into three sections, pre-S1, pre-S2, and S. Because of the multiple start codons, polypeptides of three different sizes called large (the order from surface to the inside: pre-S1, pre-S2, and S ), middle (pre-S2, S), and small (S) are produced. There is a myristyl group, which plays an important role in infection, on the amino-terminal end of the preS1 part of the large (L) protein. In addition to that, N terminus of the L protein have virus attachment and capsid binding sites. Because of that, the N termini of half of the L protein molecules are positioned outside the membrane and the other half positioned inside the membrane.
The function of the protein coded for by gene X is not fully understood but it is associated with the development of liver cancer. It stimulates genes that promote cell growth and inactivates growth regulating molecules.
Pathogenesis
The life cycle of virus is complex. is one of a few known pararetroviruses: non-retroviruses that still use reverse transcription in their replication process. The virus gains entry into the cell by binding to NTCP on the surface and being endocytosed. Because the virus multiplies via RNA made by a host enzyme, the viral genomic DNA has to be transferred to the cell nucleus by host proteins called chaperones. The partially double-stranded, circular viral DNA is then made fully double stranded by HBV DNA polymerase, transforming the genome into covalently closed circular DNA (cccDNA). This cccDNA serves as a template for transcription of four viral mRNAs by host RNA polymerase. The largest mRNA, (which is longer than the viral genome), is used to make the new copies of the genome and to make the capsid core protein and the viral DNA polymerase. These four viral transcripts undergo additional processing and go on to form progeny virions that are released from the cell or returned to the nucleus and re-cycled to produce even more copies. The long mRNA is then transported back to the cytoplasm where the virion P protein (the DNA polymerase) synthesizes DNA via its reverse transcriptase activity.
Serotypes and genotypes
The virus is divided into four major serotypes (adr, adw, ayr, ayw) based on antigenic epitopes presented on its envelope proteins, and into eight major genotypes (A–H). The genotypes have a distinct geographical distribution and are used in tracing the evolution and transmission of the virus. Differences between genotypes affect the disease severity, course and likelihood of complications, and response to treatment and possibly vaccination. There are two other genotypes I and J but they are not universally accepted as of 2015. The diversity of genotypes is not shown equally in the world. For example, A, D, and E genotypes have been seen in Africa prevalently while B and C genotypes are observed in Asia as widespread.
Genotypes differ by at least 8% of their sequence and were first reported in 1988 when six were initially described (A–F). Two further types have since been described (G and H). Most genotypes are now divided into subgenotypes with distinct properties.
Mechanisms
virus primarily interferes with the functions of the liver by replicating in hepatocytes. A functional receptor is NTCP. There is evidence that the receptor in the closely related duck hepatitis B virus is carboxypeptidase D. The virions bind to the host cell via the preS domain of the viral surface antigen and are subsequently internalized by endocytosis. HBV-preS-specific receptors are expressed primarily on hepatocytes; however, viral DNA and proteins have also been detected in extrahepatic sites, suggesting that cellular receptors for HBV may also exist on extrahepatic cells.
During HBV infection, the host immune response causes both hepatocellular damage and viral clearance. Although the innate immune response does not play a significant role in these processes, the adaptive immune response, in particular virus-specific cytotoxic T lymphocytes(CTLs), contributes to most of the liver injury associated with HBV infection. CTLs eliminate HBV infection by killing infected cells and producing antiviral cytokines, which are then used to purge HBV from viable hepatocytes. Although liver damage is initiated and mediated by the CTLs, antigen-nonspecific inflammatory cells can worsen CTL-induced immunopathology, and platelets activated at the site of infection may facilitate the accumulation of CTLs in the liver.
Diagnosis
The tests, called assays, for detection of virus infection involve serum or blood tests that detect either viral antigens (proteins produced by the virus) or antibodies produced by the host. Interpretation of these assays is complex.
The surface antigen (HBsAg) is most frequently used to screen for the presence of this infection. It is the first detectable viral antigen to appear during infection. However, early in an infection, this antigen may not be present and it may be undetectable later in the infection as it is being cleared by the host. The infectious virion contains an inner "core particle" enclosing viral genome. The icosahedral core particle is made of 180 or 240 copies of the core protein, alternatively known as core antigen, or HBcAg. During this 'window' in which the host remains infected but is successfully clearing the virus, IgM antibodies specific to the core antigen (anti-HBc IgM) may be the only serological evidence of disease. Therefore, most diagnostic panels contain HBsAg and total anti-HBc (both IgM and IgG).
Shortly after the appearance of the HBsAg, another antigen called e antigen (HBeAg) will appear. Traditionally, the presence of HBeAg in a host's serum is associated with much higher rates of viral replication and enhanced infectivity; however, variants of the virus do not produce the 'e' antigen, so this rule does not always hold true. During the natural course of an infection, the HBeAg may be cleared, and antibodies to the 'e' antigen (anti-HBe) will arise immediately afterwards. This conversion is usually associated with a dramatic decline in viral replication.
If the host is able to clear the infection, eventually the HBsAg will become undetectable and will be followed by IgG antibodies to the surface antigen and core antigen (anti-HBs and anti HBc IgG). The time between the removal of the HBsAg and the appearance of anti-HBs is called the window period. A person negative for HBsAg but positive for anti-HBs either has cleared an infection or has been vaccinated previously.
Individuals who remain HBsAg positive for at least six months are considered to be carriers. Carriers of the virus may have chronic hepatitis B, which would be reflected by elevated serum alanine aminotransferase (ALT) levels and inflammation of the liver, if they are in the immune clearance phase of chronic infection. Carriers who have seroconverted to HBeAg negative status, in particular those who acquired the infection as adults, have very little viral multiplication and hence may be at little risk of long-term complications or of transmitting infection to others. However, it is possible for individuals to enter an "immune escape" with HBeAg-negative hepatitis.
PCR tests have been developed to detect and measure the amount of HBV DNA, called the viral load, in clinical specimens. These tests are used to assess a person's infection status and to monitor treatment. Individuals with high viral loads, characteristically have ground glass hepatocytes on biopsy.
Prevention
Vaccine
Vaccines for the prevention of hepatitis B have been routinely recommended for babies since 1991 in the United States. The first dose is generally recommended within a day of birth. The hepatitis B vaccine was the first vaccine capable of preventing cancer, specifically liver cancer.
Most vaccines are given in three doses over a course of days. A protective response to the vaccine is defined as an anti-HBs antibody concentration of at least 10 mIU/ml in the recipient's serum. The vaccine is more effective in children and 95 percent of those vaccinated have protective levels of antibody. This drops to around 90% at 40 years of age and to around 75 percent in those over 60 years. The protection afforded by vaccination is long lasting even after antibody levels fall below 10 mIU/ml. For newborns of HBsAg-positive mothers: hepatitis B vaccine alone, hepatitis B immunoglobulin alone, or the combination of vaccine plus hepatitis B immunoglobulin, all prevent hepatitis B occurrence. Furthermore, the combination of vaccine plus hepatitis B immunoglobulin is superior to vaccine alone. This combination prevents HBV transmission around the time of birth in 86% to 99% of cases.
Tenofovir given in the second or third trimester can reduce the risk of mother to child transmission by 77% when combined with hepatitis B immunoglobulin and the hepatitis B vaccine, especially for pregnant women with high hepatitis B virus DNA levels. However, there is not sufficient evidence that the administration of hepatitis B immunoglobulin alone during pregnancy, might reduce transmission rates to the newborn infant. No randomized control trial has been conducted to assess the effects of hepatitis B vaccine during pregnancy for preventing infant infection.
All those with a risk of exposure to body fluids such as blood should be vaccinated, if not already. Testing to verify effective immunization is recommended and further doses of vaccine are given to those who are not sufficiently immunized.
In 10- to 22-year follow-up studies there were no cases of hepatitis B among those with a normal immune system who were vaccinated. Only rare chronic infections have been documented. Vaccination is particularly recommended for high risk groups including: health workers, people with chronic kidney failure, and men who have sex with men.
Both types of the hepatitis B vaccine, the plasma-derived vaccine (PDV) and the recombinant vaccine (RV) are of similar effectiveness in preventing infection in both healthcare workers and chronic kidney failure groups. One difference was noticed among the health worker group: the RV intramuscular route was significantly more effective compared with the RV intradermal route of administration.
Other
In assisted reproductive technology, sperm washing is not necessary for males with hepatitis B to prevent transmission, unless the female partner has not been effectively vaccinated. In females with hepatitis B, the risk of transmission from mother to child with IVF is no different from the risk in spontaneous conception.
Those at high risk of infection should be tested as there is effective treatment for those who have the disease. Groups that screening is recommended for include those who have not been vaccinated and one of the following: people from areas of the world where hepatitis B occurs in more than 2%, those with HIV, intravenous drug users, men who have sex with men, and those who live with someone with hepatitis B. Screening during pregnancy is recommended in the United States.
Treatment
Acute infection does not usually require treatment and most adults clear the infection spontaneously. Early antiviral treatment may be required in fewer than 1% of people, whose infection takes a very aggressive course (fulminant hepatitis) or who are immunocompromised. On the other hand, treatment of chronic infection may be necessary to reduce the risk of cirrhosis and liver cancer. Chronically infected individuals with persistently elevated serum alanine aminotransferase, a marker of liver damage, and HBV DNA levels are candidates for therapy. Treatment lasts from six months to a year, depending on medication and genotype. Treatment duration when medication is taken by mouth, however, is more variable and usually longer than one year.
Although none of the available medications can clear the infection, they can stop the virus from replicating, thus minimizing liver damage. As of 2024, there are seven medications licensed for the treatment of infection in the United States. These include antiviral medications lamivudine, adefovir, tenofovir disoproxil, tenofovir alafenamide, telbivudine, and entecavir, and the two immune system modulators interferon alpha-2a and PEGylated interferon alpha-2a. In 2015, the World Health Organization recommended tenofovir or entecavir as first-line agents. Those with current cirrhosis are in most need of treatment.
The use of interferon, which requires injections daily or thrice weekly, has been supplanted by long-acting PEGylated interferon, which is injected only once weekly. However, some individuals are much more likely to respond than others, and this might be because of the genotype of the infecting virus or the person's heredity. The treatment reduces viral replication in the liver, thereby reducing the viral load (the amount of virus particles as measured in the blood). Response to treatment differs between the genotypes. Interferon treatment may produce an e antigen seroconversion rate of 37% in genotype A but only a 6% seroconversion in type D. Genotype B has similar seroconversion rates to type A while type C seroconverts only in 15% of cases. Sustained e antigen loss after treatment is ~45% in types A and B but only 25–30% in types C and D.
It seems unlikely that the disease will be eliminated by 2030, the goal set in 2016 by WHO. However, progress is being made in developing therapeutic treatments. In 2010, the Hepatitis B Foundation reported that 3 preclinical and 11 clinical-stage drugs were under development, based on largely similar mechanisms. In 2020, they reported that there were 17 preclinical- and 32 clinical-stage drugs under development, using diverse mechanisms.
Prognosis
virus infection may be either acute (self-limiting) or chronic (long-standing). Persons with self-limiting infection clear the infection spontaneously within weeks to months.
Children are less likely than adults to clear the infection. More than 95% of people who become infected as adults or older children will stage a full recovery and develop protective immunity to the virus. However, this drops to 30% for younger children, and only 5% of newborns that acquire the infection from their mother at birth will clear the infection. This population has a 40% lifetime risk of death from cirrhosis or hepatocellular carcinoma. Of those infected between the age of one to six, 70% will clear the infection.
Hepatitis D (HDV) can occur only with a concomitant infection, because HDV uses the HBV surface antigen to form a capsid. Co-infection with hepatitis D increases the risk of liver cirrhosis and liver cancer. Polyarteritis nodosa is more common in people with infection.
Cirrhosis
A number of different tests are available to determine the degree of cirrhosis present. Transient elastography (FibroScan) is the test of choice, but it is expensive. Aspartate aminotransferase to platelet ratio index may be used when cost is an issue.
Reactivation
virus DNA remains in the body after infection, and in some people, including those that do not have detectable HBsAg, the disease recurs. Although rare, reactivation is seen most often following alcohol or drug use, or in people with impaired immunity. HBV goes through cycles of replication and non-replication. Approximately 50% of overt carriers experience acute reactivation. Males with baseline ALT of 200 UL/L are three times more likely to develop a reactivation than people with lower levels. Although reactivation can occur spontaneously, people who undergo chemotherapy have a higher risk. Immunosuppressive drugs favor increased HBV replication while inhibiting cytotoxic T cell function in the liver. The risk of reactivation varies depending on the serological profile; those with detectable HBsAg in their blood are at the greatest risk, but those with only antibodies to the core antigen are also at risk. The presence of antibodies to the surface antigen, which are considered to be a marker of immunity, does not preclude reactivation. Treatment with prophylactic antiviral drugs can prevent the serious morbidity associated with HBV disease reactivation.
Epidemiology
Approximately 254 million people had chronic HBV infection as of 2022. Another 1.2 million cases of acute HBV infection also occurred that year. Regional prevalences across the globe range from around 7.5% in Africa to 0.5% in the Americas.
The primary method of HBV transmission and the prevalence of chronic HBV infection in specific regions often correspond with one another. In populations where HBV infection rates are 8% or higher, which are classified as high prevalence, vertical transmission (usually occurring during birth) is most common, though rates of early childhood transmission can also be significant among these populations. In 2021, 19 African countries had infection rates ranging between 8-19%, placing them in the high prevalence category. High prevalence of HBV also exists in Mongolia.
In moderate prevalence areas where 2–7% of the population is chronically infected, the disease is predominantly spread horizontally, often among children, but also vertically. China's HBV infection rate is at the higher end of the moderate prevalence classification with an infection rate of 6.89% as of 2019. HBV prevalence in India is also moderate, with studies placing India's infection rate between 2-4%.
Countries with low HBV prevalence include Australia (0.9%), those in the WHO European Region (which average 1.5%), and most countries in North and South America (which average 0.28%).
In the United States, an estimated 0.26% of the population was living with HBV infection as of 2018.
History
Findings of HBV DNA in ancient human remains have shown that HBV has infected humans for at least ten millennia, both in Eurasia and in the Americas. This disproved the belief that hepatitis B originated in the New World and spread to Europe around 16th century. Hepatitis B virus subgenotype C4 is exclusively present in Australian aborigines, suggesting an ancient origin as much as 50,000 years old. However, analyses of ancient HBV genomes suggested that the most recent common ancestor of all known human HBV strains was dated to between 20,000 and 12,000 years ago, pointing to a more recent origin for all HBV genotypes. The evolution of HBV in humans was shown to reflect known events of human history such as the first peopling of the Americas during the late Pleistocene and the Neolithic transition in Europe. Ancient DNA studies have also showed that some ancient hepatitis viral strains still infect humans, while other strains became extinct.
The earliest record of an epidemic caused by virus was made by Lurman in 1885. An outbreak of smallpox occurred in Bremen in 1883 and 1,289 shipyard employees were vaccinated with lymph from other people. After several weeks, and up to eight months later, 191 of the vaccinated workers became ill with jaundice and were diagnosed with serum hepatitis. Other employees who had been inoculated with different batches of lymph remained healthy. Lurman's paper, now regarded as a classical example of an epidemiological study, proved that contaminated lymph was the source of the outbreak. Later, numerous similar outbreaks were reported following the introduction, in 1909, of hypodermic needles that were used, and, more importantly, reused, for administering Salvarsan for the treatment of syphilis.
The largest recorded outbreak of hepatitis B was the infection of up to 330,000 American soldiers during World War II. The outbreak has been blamed on a yellow fever vaccine made with contaminated human blood serum, and after receiving the vaccinations about 50,000 soldiers developed jaundice.
The virus was not discovered until 1966 when Baruch Blumberg, then working at the National Institutes of Health (NIH), discovered the Australia antigen (later known to be surface antigen, or HBsAg) in the blood of Aboriginal Australian people. Although a virus had been suspected since the research published by Frederick MacCallum in 1947, David Dane and others discovered the virus particle in 1970 by electron microscopy. In 1971, the FDA issued its first-ever blood supply screening order to blood banks. By the early 1980s the genome of the virus had been sequenced, and the first vaccines were being tested.
Society and culture
World Hepatitis Day, observed 28 July, aims to raise global awareness of and hepatitis C and encourage prevention, diagnosis, and treatment. It has been led by the World Hepatitis Alliance since 2007 and in May 2010, it received global endorsement from the World Health Organization.
| Biology and health sciences | Viral diseases | Health |
15929223 | https://en.wikipedia.org/wiki/Pig%20farming | Pig farming | Pig farming, pork farming, or hog farming is the raising and breeding of domestic pigs as livestock, and is a branch of animal husbandry. Pigs are farmed principally for food (e.g. pork: bacon, ham, gammon) and skins.
Pigs are amenable to many different styles of farming: intensive commercial units, commercial free range enterprises, or extensive farming (being allowed to wander around a village, town or city, or tethered in a simple shelter or kept in a pen outside the owner's house). Historically, farm pigs were kept in small numbers and were closely associated with the residence of the owner, or in the same village or town. They were valued as a source of meat and fat, and for their ability to convert inedible food into meat and manure, and were often fed household food waste when kept on a homestead. Pigs have been farmed to dispose of municipal garbage on a large scale.
All these forms of pig farm are in use today, though intensive farms are by far the most popular, due to their potential to raise a large amount of pigs in a very cost-efficient manner. In developed nations, commercial farms house thousands of pigs in climate-controlled buildings. Pigs are a popular form of livestock, with more than one billion pigs butchered each year worldwide, 100 million in the United States. The majority of pigs are used for human food, but also supply skin, fat and other materials for use in clothing, ingredients for processed foods, cosmetics, and medical use.
Production and trade
Pigs are farmed in many countries, though the countries mainly consuming them are in Asia, meaning there is a significant international and even intercontinental trade in live and slaughtered pigs. Despite having the world's largest herd, China is a net importer of pigs as China consumes about 50% of global pork production. The total amount of pork consumed in China is 57 million tons (as of 2021) and pork accounted for 60 percent of total meat consumption within the country. China has been increasing its imports during its economic development; many within China's population of 1.2 billion people prioritize eating pork as their main consumption of meat, unlike other countries where most people would prioritize having poultry. In addition, since 2007, China possesses a strategic pork reserve with a government mandate to "stabilize live hog prices, prevent excessive hog price drops, which damage the interests of farmers and to ease the negative effects of the cyclical nature of hog production and market prices." In China, the government actively intervened in the pork market during periods of instability by releasing pork reserves into the market whenever hogs get too expensive in China, in order to hold down prices for consumers. Conversely when prices of pork are deemed too low and unsustainable for farmers, the reserve buys up pigs to ensure farmers remain profitable.
The largest exporters of pigs are the United States, the European Union, and Canada. As an example, more than half of Canadian production (22.8 million pigs) in 2008 was exported, going to 143 countries.
Among animals raised for their meat, pigs have a lower feed conversion ratio than cattle, which can provide an advantage in lower unit price of meat because the cost of animal feed per kilogram or pound of resultant meat is lower. However, there are also many other economic variables in meat production and distribution, so the price differential of pork and beef at the point of retail sale does not always correspond closely to the differential in feed conversion ratios. Nonetheless, the favorable ratio often tends to make pork more affordable compared to beef.
Relationship between handlers and pigs
The way in which a stockperson interacts with pigs affects animal welfare which in some circumstances can correlate with production measures. Many routine interactions can cause fear, which can result in stress and decreased production.
There are various methods of handling pigs which can be separated into those which lead to positive or negative reactions by the animals. These reactions are based on how the pigs interpret a handler's behavior.
Negative interactions
Many negative interactions with pigs arise from stock-people dealing with large numbers of pigs. Because of this, many handlers can become complacent about animal welfare and fail to ensure positive interactions with pigs. Negative interactions include overly heavy tactile interactions (slaps, punches, kicks, and bites), the use of electric goads and fast movements. It can also include killing them. However, it is not a commonly held view that death is a negative interaction. These interactions can result in fear in the animals, which can develop into stress. Overly heavy tactile interactions from the human handlers can cause increased basal cortisol levels (a "stress" hormone). Negative interactions that cause fear mean the escape reactions of the pigs can be extremely vigorous, thereby risking injury to both stock and handlers. Stress can result in immunosuppression, leading to an increased susceptibility to disease. Studies have shown that these negative handling techniques result in an overall reduction in growth rates of pigs.
"In Canada the Federal government does not regulate the treatment on farms and most provinces have animal cruelty legislation but they typically contain expectations for general agricultural practices." Due to this lack of legislation, this perpetuates to the cruel treatment of swine. "The NFACC codes of practice are developed larger by the industry and are not enforced with third party oversight."
Positive interactions
Various interactions can be considered either positive or neutral. Neutral interactions are considered positive because, in conjunction with positive interactions, they contribute to an overall non-negative relationship between the pig handler and the animal livestock. Pigs are often fearful of fast movements. When entering a pen, it is good practice for the pig handler to enter with slow and deliberate movements. These minimize fear and therefore reduce stress. Pigs are very curious animals. Allowing the pigs to approach and smell whilst patting or resting a hand on the pig's back are examples of positive behavior. Pigs also respond positively to verbal interaction. Minimizing fear of humans allow handlers to perform husbandry practices in a safer and more efficient manner. By reducing stress, stock are made more comfortable to feed when near the pig handlers, resulting in increased productivity.
Impacts on sow breeding
Hogs raised in confinement systems tend to produce 23.5 piglets per year. Between 2013 and 2016, sow death rates nearly doubled in the United States, from 5.8 to 10.2 percent. 25 to 50 percent of deaths were caused by prolapse.
Other probable causes of death include vitamin deficiency, mycotoxins in feed, high density diets or abdominal issues. Iowa's Pork Industry Center collects mortality data in collaboration with the National Pork Board to collect data from over 400,000 sows from 16 U.S. states. The farms range in size and facility types. Increasing death rates are a profit concern to the industry, so money is invested into research to find solutions.
Genetic manipulation
Pigs were originally bred to rapidly gain weight and backfat in the late 1980s. In a more fat-conscious modern day America, pigs are now being bred to have less back fat and produce more offspring, which pushes the sow's body too far and is deemed one of the causes of the current prolapse epidemic. Researchers and veterinarians are seeking ways to positively impact the health of the hogs and benefit the hog business without taking much from the economy.
Terminology
Pigs are extensively farmed, and therefore the terminology is well developed:
Pig, hog, or swine, the species as a whole, or any member of it. The singular of "swine" is the same as the plural.
Shoat (or shote), piglet, or (where the species is called "hog") pig, unweaned young pig, or any immature pig
Sucker, a pig between birth and weaning
Weaner, a young pig recently separated from the sow
Runt, an unusually small and weak piglet, often one in a litter
Boar or hog, male pig of breeding age
Barrow, male pig castrated before puberty
Stag, male pig castrated later in life (castrated after maturity)
Gilt, young female not yet mated, or not yet farrowed, or after only one litter (depending on local usage).
Sow, breeding female, or female after first or second litter
Pigs for slaughter
Suckling pig, a piglet slaughtered for its tender meat
Feeder pig, a weaned gilt or barrow weighing between and at 6 to 8 weeks of age that is sold to be finished for slaughter
Porker, market pig between and about dressed weight
Baconer, a market pig between and dressed weight. The maximum weight can vary between processors.
Grower, a pig between weaning and sale or transfer to the breeding herd, sold for slaughter or killed for rations.
Finisher, a grower pig over liveweight
Butcher hog, a pig of approximately , ready for the market. In some markets (Italy) the final weight of butcher pig is in the range. They tend to have hind legs suitable to produce cured ham
Backfatter, cull breeding pig sold for meat; usually refers specifically to a cull sow, but is sometimes used in reference to boars
Groups
Herd, a group of pigs, or all the pigs on a farm or in a region
Sounder, a small group of pigs (or wild boar) foraging in woodland
Pig parts
Trotters, the hooves of pigs (they have four hoofed toes on each foot, walking mainly on the larger central two)
Biology
In pig, pregnant
Farrowing, giving birth
Hogging, a sow when on heat (during estrus)
Housing
Sty, a small pig-house, usually with an outdoor run or a pig confinement
Pig-shed, a larger pig-house
Ark, a low semi circular field-shelter for pigs
Curtain-barn, a long, open building with curtains on the long sides of the barn. This increases ventilation on hot, humid summer days
Environmental and health impacts
Feces and waste often spread to surrounding neighborhoods, polluting air and water with toxic waste particles. Waste from swine on these farms carry a host of pathogens and bacteria as well as heavy metals. These toxins can leach down through the soil into groundwater, polluting local drinking water supplies. Pathogens can also become airborne, polluting the air and harming individuals when ingested. Contents from waste have been shown to cause detrimental health implications, as well as harmful algal blooms in surrounding bodies of water. Due to Concentrated Animal Feed Operations (CAFOs), those who live in the surrounding areas of pig farms tend to experience health complications. Symptoms included headaches, nausea, and weakness due to the fumes that are emitted from these farms. Those who work directly inside these farms often experience these symptoms more intensely. Typically, workers of these farms experience respiratory issues such as wheezing, coughing, and tightness of the chest as well as eye and nasal irritation. This is in part due to the air quality being poor because of the air particles being contaminated with hog feces.
Little to no regulation has been written by the EPA and federal legislators surrounding CAFOs to protect the welfare of both the environment and humans from their impacts. The only permit required by federal law on wastewater runoff by CAFOs is the National Pollutant Discharge Elimination System (NPDES) permit. NPDES are authorized under the Clean Water Act and aim to reduce dumping of pollutants in water systems. However, one of the most detrimental waste management practices used at swine farms, manure lagoons, have little to no regulations surrounding waste management, as they are not connected to a moving water source and therefore is not seen as an imminent threat to human or environmental health.
Occupational hazards
Common occupational hazards faced by pig farmers include but are not limited to exposure to toxic gases and particulate matter. The Occupational Safety and Health Administration or OSHA sets health and safety standards for hazardous substances in the workplace called permissible exposure limits or PELs. Specific PELs exist for toxic gases and particulate matter and these standards are legally enforced by OSHA to ensure that the safety and health of workers are protected.
Toxic Gas and Particulate Matter Exposure
Toxic gases including hydrogen sulfide, ammonia, methane and carbon dioxide are produced as a result of the decomposition of pig feces and these gases become highly concentrated in enclosed spaces of pig barns which can be hazardous to health when inhaled. Carbon monoxide is another commonly associated toxic gas that can accumulate in pig barns as a result of the trapping of combustion byproducts such as malfunctioning furnaces or gas heat sources in the absence of adequate ventilation.
Hydrogen sulfide gas has a foul, "rotten eggs" smell at low concentrations but paralyzes the olfactory nerve at higher concentrations so that no smell is sensed. Exposure to high levels, well beyond the OSHA PEL, of hydrogen sulfide can cause fatal respiratory paralysis. The common source of hydrogen sulfide are covered manure pits below the pig barns that act as feces reservoirs. These manure pits require regular emptying and during this process, high levels of hydrogen sulfide is released and seeps into pig barns. Pig barns must be void of any human or animal inhabitants during this emptying process and require a several hour "waiting period" until occupants can safely reenter the barn.
Ammonia gas has a strong odor that can be smelled at low levels, below the OSA PEL, but does not have any negative health effects. At higher levels, ammonia is irritating to the body's mucous membranes such as the eyes, nose, mouth, throat and lungs. Particulate matter in pig barns often absorbs ammonia as it floats through the air. These particles are then inhaled and increase the irritating effect of ammonia.
Methane and carbon dioxide are combustible gases meaning that they can burn, catch fire or explode easily. They are also known as chemical asphyxiants and at high levels can cause suffocation by displacing oxygen from the air.
Particulate matter is produced when small fragments of pig hair or skin, dried feces, or feed can detach and become suspended in the air in pig barns. The increased concentration of particulate matter in the air, especially in confined spaces, can lead to respiratory tract irritation and other health effects when inhaled. Bacteria and viruses, such as influenza, can travel through the air on particulate matter and increase the risk of transmission of disease.
OSHA requires that toxic gas and particulate matter be measured at least twice yearly preferably in the autumn months and again in the winter when natural ventilation is the most reduced. Workers are also advised to wear N-95 respirators and eye protection when inside of pig barns to prevent the inhalation of toxic gases and particulate matter as well as irritation to the of eyes.
Geopolitical issues
As with other commodities, pork presents challenges in the politics of international trade as national interests compete and seek economic modus vivendi. Changes to policy can upset the existing balances, prompting economic anxiety. For example, in 2020, the hog farming sector in Taiwan was upset by a decision to allow imports from the United States without labeling of ractopamine use. Farmers' views varied on how negative the effects might be. Issues of pride and degree of autarky also figure into such debates; people understandably wonder whether trade competition changes will deeply damage domestic production capability, while accurate quantitative answers are often difficult to find amid the mass of debate.
Drugs
Growth promoters
Ractopamine
Most pigs in the US receive ractopamine which promotes muscle instead of fat, quicker weight gain, and reduced costs and pollutants in the environment. Such pigs consume less feed to reach finishing weight and produce less manure. Ractopamine has not been approved for use by the European Union, China, Russia, and several other countries.
Colistin
China once used colistin (an antibiotic) as growth promoter (subtherapeutic antibiotic use) but discovered a colistin-resistant form of E. coli bacteria in a pig from a Shanghai farm in 2013. Investigations then led to the identification of "a gene called MCR-1 that allowed bacteria to survive colistin treatment in animals and humans." In 2016, these findings led China to ban colistin as growth promoter.
Antibiotics
A systematic review found that penicillins and tetracyclines were the most commonly used antibiotics in pigs.
Parasites
Toxoplasmosis is a constant pressure on pig farming. Worldwide, the percentage of pigs harboring viable Toxoplasma gondii parasites has been measured to be 3% to 71.43%. Surveys of seroprevalence (T. gondii antibodies in blood) are more common, and such measurements are indicative of the high relative seroprevalence in pigs across the world. Neonatal piglets have been found to suffer the entire range of severity, including progression to stillbirth. This was especially demonstrated in the foundational Thiptara et al. 2006, reporting a litter birth of three stillborns and six live in Thailand. This observation has been relevant not only to that country but to toxoplasmosis control in porciculture around the world.
Hygiene
Excessively hygienic raising conditions were found to prevent proper gut microbiota development by Schmidt et al. 2011. Moore et al. 1995 describes the pathology of Cryptosporidium infection, a common difficulty in piglet production.
In an attempt to curb diseases such as African swine fever, a number of Chinese companies have built condominium-style mega complexes multiple stories high to house thousands of pigs. The buildings have been dubbed "hog hotels" and come with strict protocols and advanced cleaning, veterinary, and disposal systems. However, doubt has been raised by policy specialists and animal scientists over the facilities' efficacy in preventing outbreaks. The welfare of the animals has also been a source of concern, and it has been suggested that the poor welfare of the pigs may cause a decline in their immunity.
| Technology | Animal husbandry | null |
5457188 | https://en.wikipedia.org/wiki/Parasitoid%20wasp | Parasitoid wasp | Parasitoid wasps are a large group of hymenopteran superfamilies, with all but the wood wasps (Orussoidea) being in the wasp-waisted Apocrita. As parasitoids, they lay their eggs on or in the bodies of other arthropods, sooner or later causing the death of these hosts. Different species specialise in hosts from different insect orders, most often Lepidoptera, though some select beetles, flies, or bugs; the spider wasps (Pompilidae) exclusively attack spiders.
Parasitoid wasp species differ in which host life-stage they attack: eggs, larvae, pupae, or adults. They mainly follow one of two major strategies within parasitism: either they are endoparasitic, developing inside the host, and koinobiont, allowing the host to continue to feed, develop, and moult; or they are ectoparasitic, developing outside the host, and idiobiont, paralysing the host immediately. Some endoparasitic wasps of the superfamily Ichneumonoidea have a mutualistic relationship with polydnaviruses, the viruses suppressing the host's immune defenses.
Parasitoidism evolved only once in the Hymenoptera, during the Permian, leading to a single clade called Euhymenoptera, but the parasitic lifestyle has secondarily been lost several times including among the ants, bees, and vespid wasps. As a result, the order Hymenoptera contains many families of parasitoids, intermixed with non-parasitoid groups. The parasitoid wasps include some very large groups, some estimates giving the Chalcidoidea as many as 500,000 species, the Ichneumonidae 100,000 species, and the Braconidae up to 50,000 species.
Host insects have evolved a range of defences against parasitoid wasps, including hiding, wriggling, and camouflage markings.
Many parasitoid wasps are considered beneficial to humans because they naturally control agricultural pests. Some are applied commercially in biological pest control, starting in the 1920s with Encarsia formosa to control whitefly in greenhouses. Historically, parasitoidism in wasps influenced the thinking of Charles Darwin.
Parasitoidism
Parasitoid wasps range from some of the smallest species of insects to wasps about an inch long. Most females have a long, sharp ovipositor at the tip of the abdomen, sometimes lacking venom glands, and almost never modified into a sting.
Parasitoids can be classified in a variety of ways. They can live within their host's body as endoparasitoids, or feed on it from outside as ectoparasitoids: both strategies are found among the wasps. Parasitoids can also be divided according to their effect on their hosts. Idiobionts prevent further development of the host after initially immobilizing it, while koinobionts allow the host to continue its development while they are feeding upon it; and again, both types are seen in parasitoidal wasps. Most ectoparasitoid wasps are idiobiont, as the host could damage or dislodge the external parasitoid if allowed to move or moult. Most endoparasitoid wasps are koinobionts, giving them the advantage of a host that continues to grow larger and remains able to avoid predators.
Hosts
Many parasitoid wasps use larval Lepidoptera as hosts, but some groups parasitize different host life stages (egg, larva or nymph, pupa, adult) of nearly all other orders of insects, especially Coleoptera, Diptera, Hemiptera and other Hymenoptera. Some attack arthropods other than insects: for instance, the Pompilidae specialise in catching spiders: these are quick and dangerous prey, often as large as the wasp itself, but the spider wasp is quicker, swiftly stinging her prey to immobilise it. Adult female wasps of most species oviposit into their hosts' bodies or eggs.
More rarely, parasitoid wasps may use plant seeds as hosts, such as Torymus druparum.
Some also inject a mix of secretory products that paralyse the host or protect the egg from the host's immune system; these include polydnaviruses, ovarian proteins, and venom. If a polydnavirus is included, it infects the nuclei of host hemocytes and other cells, causing symptoms that benefit the parasite.
Host size is important for the development of the parasitoid, as the host is its entire food supply until it emerges as an adult; small hosts often produce smaller parasitoids. Some species preferentially lay female eggs in larger hosts and male eggs in smaller hosts, as the reproductive capabilities of males are limited less severely by smaller adult body size.
Some parasitoid wasps mark the host with chemical signals to show that an egg has been laid there. This may both deter rivals from ovipositing, and signal to itself that no further egg is needed in that host, effectively reducing the chances that offspring will have to compete for food and increasing the offspring's survival.
Life cycle
On or inside the host the parasitoid egg hatches into a larva or two or more larvae (polyembryony). Endoparasitoid eggs can absorb fluids from the host body and grow several times in size from when they were first laid before hatching. The first instar larvae are often highly mobile and may have strong mandibles or other structures to compete with other parasitoid larvae. The following instars are generally more grub-like. Parasitoid larvae have incomplete digestive systems with no rear opening. This prevents the hosts from being contaminated by their wastes. The larva feeds on the host's tissues until ready to pupate; by then the host is generally either dead or almost so. A meconium, or the accumulated wastes from the larva is cast out as the larva transitions to a prepupa. Depending on its species, the parasitoid then may eat its way out of the host or remain in the more or less empty skin. In either case it then generally spins a cocoon and pupates. As adults, parasitoid wasps feed primarily on nectar from flowers. Females of some species will also drink hemolymph from hosts to gain additional nutrients for egg production.
Mutualism with polydnavirus
Polydnaviruses are a unique group of insect viruses that have a mutualistic relationship with some parasitic wasps. The polydnavirus replicates in the oviducts of an adult female parasitoid wasp. The wasp benefits from this relationship because the virus provides protection for the parasitic larvae inside the host, (i) by weakening the host's immune system and (ii) by altering the host's cells to be more beneficial to the parasite. The relationship between these viruses and the wasp is obligatory in the sense that all individuals are infected with the viruses; the virus has been incorporated in the wasp's genome and is inherited.
Host defenses
The hosts of parasitoids have developed several levels of defence. Many hosts try to hide from the parasitoids in inaccessible habitats. They may also get rid of their frass (body wastes) and avoid plants that they have chewed on as both can signal their presence to parasitoids hunting for hosts. The egg shells and cuticles of the potential hosts are thickened to prevent the parasitoid from penetrating them. Hosts may use behavioral evasion when they encounter an egg laying female parasitoid, like dropping off the plant they are on, twisting and thrashing so as to dislodge or kill the female and even regurgitating onto the wasp to entangle it. The wriggling can sometimes help by causing the wasp to "miss" laying the egg on the host and instead place it nearby. Wriggling of pupae can cause the wasp to lose its grip on the smooth hard pupa or get trapped in the silk strands. Some caterpillars even bite the female wasps that approach them. Some insects secrete poisonous compounds that kill or drive away the parasitoid. Ants that are in a symbiotic relationship with caterpillars, aphids or scale insects may protect them from attack by wasps.
Parasitoid wasps are vulnerable to hyperparasitoid wasps. Some parasitoid wasps change the behavior of the infected host, causing them to build a silk web around the pupae of the wasps after they emerge from its body to protect them from hyperparasitoids.
Hosts can kill endoparasitoids by sticking haemocytes to the egg or larva in a process called encapsulation. In aphids, the presence of a particular species of γ-3 Pseudomonadota makes the aphid relatively immune to their parasitoid wasps by killing many of the eggs. As the parasitoid's survival depends on its ability to evade the host's immune response, some parasitoid wasps have developed the counterstrategy of laying more eggs in aphids that have the endosymbiont, so that at least one of them may hatch and parasitize the aphid.
Certain caterpillars eat plants that are toxic to both themselves and the parasite to cure themselves. Drosophila melanogaster larvae also self-medicate with ethanol to treat parasitism. D. melanogaster females lay their eggs in food containing toxic amounts of alcohol if they detect parasitoid wasps nearby. The alcohol protects them from the wasps, at the cost of retarding their own growth.
Evolution and taxonomy
Evolution
Based on genetic and fossil analysis, parasitoidism has evolved only once in the Hymenoptera, during the Permian, leading to a single clade. All parasitoid wasps are descended from this lineage. The narrow-waisted Apocrita emerged during the Jurassic. The Aculeata, which includes bees, ants, and parasitoid spider wasps, evolved from within the Apocrita; it contains many families of parasitoids, though not the Ichneumonoidea, Cynipoidea, and Chalcidoidea. The Hymenoptera, Apocrita, and Aculeata are all clades, but since each of these contains non-parasitic species, the parasitoid wasps, formerly known as the Parasitica, do not form a clade on their own. The common ancestor in which parasitoidism evolved lived approximately 247 million years ago and was previously believed to be an ectoparasitoid wood wasp that fed on wood-boring beetle larvae. Species similar in lifestyle and morphology to this ancestor still exist in the Ichneumonoidea. However, recent molecular and morphological analysis suggests this ancestor was endophagous, meaning it fed from within its host. A significant radiation of species in the Hymenoptera occurred shortly after the evolution of parasitoidy in the order and is thought to have been a result of it. The evolution of a wasp waist, a constriction in the abdomen of the Apocrita, contributed to rapid diversification as it increased maneuverability of the ovipositor, the organ off the rear segment of the abdomen used to lay eggs.
The phylogenetic tree gives a condensed overview of the positions of parasitoidal groups (boldface), amongst groups (italics) like the Vespidae which have secondarily abandoned the parasitoid habit. The approximate numbers of species estimated to be in these groups, often much larger than the number so far described, is shown in parentheses, with estimates for the most populous also shown in boldface, like "(150,000)". Not all species in these groups are parasitoidal: for example, some Cynipoidea are phytophagous.
Taxonomy
The parasitoid wasps are paraphyletic since the ants, bees, and non-parasitic wasps such as the Vespidae are not included, and there are many members of mainly parasitoidal families which are not themselves parasitic. Listed are Hymenopteran families where most members have a parasitoid lifestyle.
Symphyta:
Orussidae
Apocrita:
Scolebythidae
Bethylidae
Chrysididae
Sclerogibbidae
Dryinidae
Embolemidae
Tiphiidae
Thynnidae
Sapygidae
Mutillidae
Bradynobaenidae
Chyphotidae
Sierolomorphidae
Braconidae
Ichneumonidae
Pompilidae
Rhopalosomatidae
Aulacidae
Evaniidae
Gasteruptiidae
Stephanidae
Megalyridae
Trigonalidae
Ibaliidae
Liopteridae
Figitidae
Austroniidae
Diapriidae
Heloridae
Monomachidae
Pelecinidae
Peradeniidae
Proctotrupidae
Roproniidae
Vanhorniidae
Platygastridae
Scelionidae
Megaspilidae
Ceraphronidae
Mymarommatidae
Chalcidoidea (19 families)
Ampulicidae
Interactions with humans
Biological pest control
Parasitoid wasps are considered beneficial as they naturally control the population of many pest insects. They are widely used commercially (alongside other parasitoids such as tachinid flies) for biological pest control, for which the most important groups are the ichneumonid wasps, which prey mainly on caterpillars of butterflies and moths; braconid wasps, which attack caterpillars and a wide range of other insects including greenfly; chalcidoid wasps, which parasitise eggs and larvae of greenfly, whitefly, cabbage caterpillars, and scale insects.
One of the first parasitoid wasps to enter commercial use was Encarsia formosa, an endoparasitic aphelinid. It has been used to control whitefly in greenhouses since the 1920s. Use of the insect fell almost to nothing, replaced by chemical pesticides by the 1940s. Since the 1970s, usage has revived, with renewed usage in Europe and Russia. In some countries, such as New Zealand, it is the primary biological control agent used to control greenhouse whiteflies, particularly on crops such as tomato, a particularly difficult plant for predators to establish on.
Commercially, there are two types of rearing systems: short-term seasonal daily output with high production of parasitoids per day, and long-term year-round low daily output with a range in production of 4–1000 million female parasitoids per week, to meet demand for suitable parasitoids for different crops.
In culture
Parasitoid wasps influenced the thinking of Charles Darwin. In an 1860 letter to the American naturalist Asa Gray, Darwin wrote: "I cannot persuade myself that a beneficent and omnipotent God would have designedly created parasitic wasps with the express intention of their feeding within the living bodies of Caterpillars." The palaeontologist Donald Prothero notes that religiously-minded people of the Victorian era, including Darwin, were horrified by this instance of evident cruelty in nature, particularly noticeable in the Ichneumonidae.
| Biology and health sciences | Hymenoptera | Animals |
2155942 | https://en.wikipedia.org/wiki/Harvester%20ant | Harvester ant | Harvester ant is a common name for any of the species or genera of ants that collect seeds (called seed predation), or mushrooms as in the case of Euprenolepis procera, which are stored in the nest in communal chambers called granaries. They are also referred to as agricultural ants. Seed harvesting by some desert ants is an adaptation to the lack of typical ant resources such as prey or honeydew from hemipterans. Harvester ants increase seed dispersal and protection, and provide nutrients that increase seedling survival of the desert plants. In addition, ants provide soil aeration through the creation of galleries and chambers, mix deep and upper layers of soil, and incorporate organic refuse into the soil.
Seed dispersal
Ants may play an important role in the dynamics of plant communities by acting either as seed dispersal agents or as seed predators, or both. During the day, these ants search the savannas for vegetation and plant seeds, and carry them along back to their nest. The two main mechanisms through which ants disperse seeds are myrmecochory, or seed dispersal mediated by the elaiosome, i.e., a lipid-rich seed appendage that mainly attracts non-granivorous ants and provides rewards for seed dispersal, and diszoochory, or seed dispersal performed by seed-harvesting ants that is not mediated by any particular seed structure. While the former has traditionally been recognized mainly as a mutualism, the latter is usually perceived as an antagonism.
Foraging behavior
Harvester ants foraging in hot, dry conditions lose water, but obtain water from metabolizing fats in the seeds they eat. Positive feedback on foraging activity, from returning foragers with food, allows the colony to regulate its foraging activity according to the current costs of desiccation and the benefits based on current food availability.
In many harvester ant species, foraging behavior is influenced by the weather. For example, in the ant Messor andrei, recruitment to food bait is higher in more humid conditions. Both humidity and food availability are affected by day-to-day changes in weather conditions. Food is distributed by wind and flooding and rain uncover seeds in the top layer of the soil. In Pogonomyrmex barbatus, daily changes in conditions such as humidity and food availability produce strong daily trends in the foraging activity of all colonies.
Colonies may vary in the relation between humidity and foraging activity. Colonies differ consistently from year to year in how often they forage at all and most colonies forage on days with high humidity and high food availability, such as those just after a rain when flooding has exposed a layer of seeds in the soil. Few colonies forage on very dry days. Colonies also differ in how likely they are to adjust the rate of outgoing foragers to the rate of forager return. While all colonies tend to adjust outgoing foraging rate closely when conditions are good, only some colonies do so in poor conditions.
Sting
Harvester ants, for their size, have a rather potent venom. They inject it into their victim via sting by biting down and following up with a rapid sting from their abdomen. This causes 4-8 hours of sharp pain with effects similar to neurotoxicity such as piloerection and localized swelling around the area of the sting.
Species and genera
Aphaenogaster - about 200 species
Novomessor, seed harvesters
Novomessor cockerelli
Euprenolepis - eight species
Euprenolepis procera, nomadic mushroom-harvesters, a previously unknown lifestyle among ants
Messor, seed-harvesters
Pheidole, seed-harvesters
Pogonomyrmex, seed-harvesters
Pogonomyrmex barbatus
Pogonomyrmex maricopa, a venomous species found in Arizona, USA
Pogonomyrmex occidentalis, seed-harvesters
Carebara
Carebara diversa, seed harvesters
| Biology and health sciences | Hymenoptera | Animals |
2156142 | https://en.wikipedia.org/wiki/Australorp | Australorp | The Australorp is an Australian breed of dual-purpose utility chicken. It derives from the British Black Orpington, and was selectively bred for egg-laying performance; some hens lay more than 300 eggs per year. It achieved world-wide popularity in the 1920s after the breed broke numerous world records for number of eggs laid and has been a popular breed in the western world since. It is one of eight poultry breeds created in Australia and recognised by the Australian Poultry Standards. The original plumage colour is black, which is the only colour recognised in the United States of America, but blue and white are also recognised in Australia and the Poultry Club South Africa recognises buff, splash, wheaten laced and golden in addition.
History
The original stock used in the development of the Australorp was imported to Australia from England out of the Black Orpington yards of William Cook and Joseph Partington in the period from 1890 to the early 1900s with Rhode Island Red. Local breeders used this stock together with judicious out-crossings of Minorca, White Leghorn and Langshan blood to improve the utility features of the imported Orpingtons. There is even a report of some Plymouth Rock blood also being used. The emphasis of the early breeders was on utility features. At this time, the resulting birds were known as Australian Black Orpingtons (Austral-orp).
The origin of the name "Australorp" seems to be shrouded in as much controversy as the attempts to obtain agreement between the States over a suitable national Standard. The earliest claim to the name was made by one of poultry fancy's institutions, Wiliam Wallace Scott, before the First World War. From 1925 Wal Scott set to work to have Australorp recognised as a breed with the Poultry Society as he developed the breed. Equally as persuasive a claim came in 1919 from Arthur Harwood who suggested that the "Australian Laying Orpingtons" be named "Australs". The letters "orp" were suggested as a suffix to denote the major breed in the fowl's development. A further overseas claim to the name came from Britain's W. Powell-Owen who drafted the British Standard for the breed in 1921 following the importation of the "Australian Utility Black Orpingtons". It is certain that the name "Australorp" was being used in the early 1920s when the breed was launched internationally. The Australorp was added to the Standard of Perfection of the American Poultry Association in 1929.
A white colour variety was bred in South Africa the 1930s, and in Australia in the 1940s. It was recognised in Australia in 2012. A bantam Australorp was bred in the early 1930s by Roy Corner and Jack Mann, and was first exhibited at a poultry show in 1934.
The Australorp is reported in all five inhabited continents, in seventeen countries of which four report a population of or more: Australia, Serbia, Slovakia and the United States. In 2024 the world-wide population was estimated at close to , with an overall conservation status of "not at risk". The risk status in Australia was listed as "at/risk/endangered" in DAD-IS, while in the United Kingdom it was listed as "priority" on the watchlist of the Rare Breeds Survival Trust. In the United States it was removed from the watchlist of the Livestock Conservancy in 2023.
Characteristics
There are both bantam and standard-sized Australorps.
Three colours are recognised in the Australian Poultry Standards: black, white and blue; the same three colours are recognised by the Entente Européenne and by the Poultry Club of Great Britain. The Poultry Club South Africa recognises a further four colours: buff, splash, wheaten laced, and golden.
Use
The egg-laying performance of Australorps attracted attention when in 1922–1923, a team of six hens set a world record by laying 1857 eggs for an average of 309.5 eggs per hen during a 365 consecutive day trial. These figures were achieved without the lighting regimens of the modern intensive shed. Such performances had importation orders flooding in from England, United States of America, South Africa, Canada and Mexico. A new record was set when a hen laid 364 eggs in 365 days. They are also known to be good nest sitters and mothers, making them one of the most popular large heritage utility breeds of chicken.
Hens lay approximately 190 light brown eggs per year, with an average weight of ; bantam hens lay some 160 per year, averaging in weight.
| Biology and health sciences | Chickens | Animals |
2158298 | https://en.wikipedia.org/wiki/Visual%20impairment | Visual impairment | Visual or vision impairment (VI or VIP) is the partial or total inability of visual perception. In the absence of treatment such as corrective eyewear, assistive devices, and medical treatment, visual impairment may cause the individual difficulties with normal daily tasks, including reading and walking. The terms low vision and blindness are often used for levels of impairment which are difficult or impossible to correct and significantly impact daily life. In addition to the various permanent conditions, fleeting temporary vision impairment, amaurosis fugax, may occur, and may indicate serious medical problems.
The most common causes of visual impairment globally are uncorrected refractive errors (43%), cataracts (33%), and glaucoma (2%). Refractive errors include near-sightedness, far-sightedness, presbyopia, and astigmatism. Cataracts are the most common cause of blindness. Other disorders that may cause visual problems include age-related macular degeneration, diabetic retinopathy, corneal clouding, childhood blindness, and a number of infections. Visual impairment can also be caused by problems in the brain due to stroke, premature birth, or trauma, among others. These cases are known as cortical visual impairment. Screening for vision problems in children may improve future vision and educational achievement. Screening adults without symptoms is of uncertain benefit. Diagnosis is by an eye exam.
The World Health Organization (WHO) estimates that 80% of visual impairment is either preventable or curable with treatment. This includes cataracts, the infections river blindness and trachoma, glaucoma, diabetic retinopathy, uncorrected refractive errors, and some cases of childhood blindness. Many people with significant visual impairment benefit from vision rehabilitation, changes in their environment, and assistive devices.
, there were 940 million people with some degree of vision loss. 246 million had low vision and 39 million were blind. The majority of people with poor vision are in the developing world and are over the age of 50 years. Rates of visual impairment have decreased since the 1990s. Visual impairments have considerable economic costs both directly due to the cost of treatment and indirectly due to decreased ability to work.
Classification
In 2010, the WHO definition for visual impairment was changed and now follows the ICD-11. The previous definition which used "best corrected visual acuity" was changed to "presenting visual acuity". This change was made as newer studies showed that best-corrected vision overlooks a larger proportion of the population who has visual impairment due to uncorrected refractive errors, and/or lack of access to medical or surgical treatment.
Distance vision impairment:
Category 0: No or mild visual impairment – presenting visual acuity better than 6/18
Category 1: Moderate visual impairment – presenting visual acuity worse than 6/18 and better than 6/60
Category 2: Severe visual impairment – presenting visual acuity worse than 6/60 and better than 3/60
Category 3: Blindness – presenting visual acuity worse than 3/60 and better than 1/60
Category 4: Blindness – presenting visual acuity worse than 1/60 with light perception
Category 5: Blindness – irreversible blindness with no light perception
Near vision impairment:
Near visual acuity worse than N6 or M 0.8 at 40 cm.
United Kingdom
Severely sight impaired
Defined as having central visual acuity of less than 3/60 with normal fields of vision, or gross visual field restriction.
Unable to see at what the normally sighted person sees at .
Sight impaired
Able to see at , but not at , what the normally sighted person sees at
Less severe visual impairment is not captured by registration data, and its prevalence is difficult to quantify
Low vision
A visual acuity of less than 6/18 but greater than 3/60.
Not eligible to drive and may have difficulty recognising faces across a street, watching television, or choosing clean, unstained, co-ordinated clothing.
In the UK, the Certificate of Vision Impairment (CVI) is used to certify people as being severely sight impaired or sight impaired. The accompanying guidance for clinical staff states: "The National Assistance Act 1948 states that a person can be certified as severely sight impaired if they are 'so blind as to be unable to perform any work for which eye sight is essential'". Certification is based on whether a person can do any work for which eyesight is essential, not just one particular job (such as their job before becoming blind).
In practice, the definition depends on individuals' visual acuity and the extent to which their field of vision is restricted. The Department of Health identifies three groups of people who may be classified as severely visually impaired.
Those below 3/60 (equivalent to 20/400 in US notation) Snellen (most people below 3/60 are severely sight impaired).
Those better than 3/60 but below 6/60 Snellen (people who have a very contracted field of vision only).
Those 6/60 Snellen or above (people in this group who have a contracted field of vision especially if the contraction is in the lower part of the field).
The Department of Health also state that a person is more likely to be classified as severely visually impaired if their eyesight has failed recently or if they are an older individual, both groups being perceived as less able to adapt to their vision loss.
United States
In the United States, any person with vision that cannot be corrected to better than 20/200 in the better eye, or who has 20 degrees (diameter) or less of visual field remaining, is considered legally blind or eligible for disability classification and possible inclusion in certain government sponsored programs.
The terms partially sighted, low vision, legally blind and totally blind are used by schools, colleges, and other educational institutions to describe students with visual impairments. They are defined as follows:
Partially sighted indicates some type of visual problem, with a need of person to receive special education in some cases.
Low vision generally refers to a severe visual impairment, not necessarily limited to distance vision. Low vision applies to all individuals with sight who are unable to read the newspaper at a normal viewing distance, even with the aid of eyeglasses or contact lenses. They use a combination of vision and other senses to learn, although they may require adaptations in lighting or the size of print, and, sometimes, braille.
Legally blind indicates that a person has less than 20/200 vision in the better eye after best correction (contact lenses or glasses), or a field of vision of less than 20 degrees in the better eye.
Totally blind students learn via braille or other non-visual media.
In 1934, the American Medical Association adopted the following definition of blindness:
The United States Congress included this definition as part of the Aid to the Blind program in the Social Security Act passed in 1935. In 1972, the Aid to the Blind program and two others combined under Title XVI of the Social Security Act to form the Supplemental Security Income program which states:
Temporary vision impairment
Vision impairment for a few seconds, or minutes, may occur due to any of a variety of causes, some serious and requiring medical attention.
Health effects
General functioning
Visual impairments may take many forms and be of varying degrees. Visual acuity alone is not always a good predictor of an individual's function. Someone with relatively good acuity (e.g., 20/40) can have difficulty with daily functioning, while someone with worse acuity (e.g., 20/200) may function reasonably well if they have low visual demands.
Best-corrected visual acuity differs from presenting visual acuity; a person with a "normal" best corrected acuity can have "poor" presenting acuity (e.g. individual who has uncorrected refractive error). Thus, measuring an individual's general functioning depends on one's situational and contextual factors, as well as access to treatment.
The American Medical Association has estimated that the loss of one eye equals 25% impairment of the visual system and 24% impairment of the whole person; total loss of vision in both eyes is considered to be 100% visual impairment and 85% impairment of the whole person.
Some people who fall into this category can use their considerable residual vision – their remaining sight – to complete daily tasks without relying on alternative methods. The role of a low vision specialist (optometrist or ophthalmologist) is to maximize the functional level of a patient's vision by optical or non-optical means. Primarily, this is by use of magnification in the form of telescopic systems for distance vision and optical or electronic magnification for near tasks.
People with significantly reduced acuity may benefit from training conducted by individuals trained in the provision of technical aids. Low vision rehabilitation professionals, some of whom are connected to an agency for the blind, can provide advice on lighting and contrast to maximize remaining vision. These professionals also have access to non-visual aids, and can instruct patients in their uses.
Mobility
Older adults with visual impairment are at an increased risk of physical inactivity, slower gait speeds, and fear of falls.
Physical activity is a useful predictor of overall well-being, and routine physical activity reduces the risk of developing chronic diseases and disability. Older adults with visual impairment (including glaucoma, age-related macular degeneration, and diabetic retinopathy) have decreased physical activity as measured with self-reports and accelerometers. The US National Health and Nutrition Examination Survey (NHANES) showed that people with corrected visual acuity of less than 20/40 spent significantly less time in moderate to vigorous physical activity. Age-related macular degeneration is also associated with a 50% decrease in physical activity–however physical activity is protective against age-related macular degeneration progression.
In terms of mobility, those with visual impairment have a slower gait speed than those without visual impairment; however, the rate of decline remains proportional with increasing age in both groups. Additionally, the visually impaired also have greater difficulty walking a quarter mile (400 m) and walking up stairs, as compared to those with normal vision.
Cognitive
Older adults with vision loss are at an increased risk of memory loss, cognitive impairment, and cognitive decline.
Social and psychological
Studies demonstrate an association between older adults with visual impairment and a poor mental health; discrimination was identified as one of the causes of this association. Older adults with visual impairment have a 1.5-fold risk of reporting perceived discrimination and of these individuals, there was a 2-fold risk of loneliness and 4-fold risk of reporting a lower quality of life. Among adults with visual impairment, the prevalence of moderate loneliness is 28.7% (18.2% in general population) and prevalence of severe loneliness is 19.7% (2.7% in general population). The risk of depression and anxiety are also increased in the visually impaired; 32.2% report depressive symptoms (12.01% in general population), and 15.61% report anxiety symptoms (10.69% in general population).
The subjects making the most use of rehabilitation instruments, who lived alone, and preserved their own mobility and occupation were the least depressed, with the lowest risk of suicide and the highest level of social integration.
Those with worsening sight and the prognosis of eventual blindness are at comparatively high risk of suicide and thus may be in need of supportive services. Many studies have demonstrated how rapid acceptance of the serious visual impairment has led to a better, more productive compliance with rehabilitation programs. Moreover, psychological distress has been reported to be at its highest when sight loss is not complete, but the prognosis is unfavorable. Therefore, early intervention is imperative for enabling successful psychological adjustment.
Associated conditions
Blindness can occur in combination with such conditions as intellectual disability, autism spectrum disorders, cerebral palsy, hearing impairments, and epilepsy. Blindness in combination with hearing loss is known as deafblindness.
It has been estimated that over half of completely blind people have non-24-hour sleep–wake disorder, a condition in which a person's circadian rhythm, normally slightly longer than 24 hours, is not entrained (synchronized) to the lightdark cycle.
Cause
The most common causes of visual impairment globally in 2010 were:
Refractive error (42%)
Cataract (33%)
Glaucoma (2%)
Age-related macular degeneration (1%)
Corneal opacification (1%)
Diabetic retinopathy (1%)
Childhood blindness
Trachoma (1%)
Undetermined (18%)
The most common causes of blindness worldwide in 2010 were:
Cataracts (51%)
Glaucoma (8%)
Age-related macular degeneration (5%)
Corneal opacification (4%)
Childhood blindness (4%)
Refractive errors (3%)
Trachoma (3%)
Diabetic retinopathy (1%)
Undetermined (21%)
About 90% of people who are visually impaired live in the developing world. Age-related macular degeneration, glaucoma, and diabetic retinopathy are the leading causes of blindness in the developed world.
Among working-age adults who are newly blind in England and Wales the most common causes in 2010 were:
Hereditary retinal disorders (20.2%)
Diabetic retinopathy (14.4%)
Optic atrophy (14.1%)
Glaucoma (5.9%)
Congenital abnormalities (5.1%)
Disorders of the visual cortex (4.1%)
Cerebrovascular disease (3.2%)
Degeneration of the macula and posterior pole (3.0%)
Myopia (2.8%)
Corneal disorders (2.6%)
Malignant neoplasms of the brain and nervous system (1.5%)
Retinal detachment (1.4%)
Cataracts
Cataracts are the greying or opacity of the crystalline lens, which can be caused in children by intrauterine infections, metabolic disorders, and genetically transmitted syndromes. Cataracts are the leading cause of child and adult blindness that doubles in prevalence with every ten years after the age of 40. Consequently, today cataracts are more common among adults than in children. That is, people face higher chances of developing cataracts as they age. Nonetheless, cataracts tend to have a greater financial and emotional toll upon children as they must undergo expensive diagnosis, long term rehabilitation, and visual assistance. Also, according to the Saudi Journal for Health Sciences, sometimes people experience irreversible amblyopia after pediatric cataract surgery because the cataracts prevented the normal maturation of vision prior to operation. Despite the great progress in treatment, cataracts remain a global problem in both economically developed and developing countries. At present, with the variant outcomes as well as the unequal access to cataract surgery, the best way to reduce the risk of developing cataracts is to avoid smoking and extensive exposure to sun light (i.e. UV-B rays).
Glaucoma
Glaucoma is an eye disease often characterized by increased pressure within the eye or intraocular pressure (IOP). Glaucoma causes visual field loss as well as severs the optic nerve. Early diagnosis and treatment of glaucoma in patients is imperative because glaucoma is triggered by non-specific levels of IOP. Also, another challenge in accurately diagnosing glaucoma is that the disease has four causes: 1) inflammatory ocular hypertension syndrome (IOHS); 2) severe uveitic angle closure; 3) corticosteroid-induced; and 4) a heterogonous mechanism associated with structural change and chronic inflammation. In addition, often pediatric glaucoma differs greatly in cause and management from the glaucoma developed by adults. Currently, the best sign of pediatric glaucoma is an IOP of 21 mm Hg or greater present within a child. One of the most common causes of pediatric glaucoma is cataract removal surgery, which leads to an incidence rate of about 12.2% among infants and 58.7% among 10-year-olds.
Infections
Childhood blindness can be caused by conditions related to pregnancy, such as congenital rubella syndrome and retinopathy of prematurity. Leprosy and onchocerciasis each blind approximately 1 million individuals in the developing world.
The number of individuals blind from trachoma has decreased in the past 10 years from 6 million to 1.3 million, putting it in seventh place on the list of causes of blindness worldwide.
Central corneal ulceration is also a significant cause of monocular blindness worldwide, accounting for an estimated 850,000 cases of corneal blindness every year in the Indian subcontinent alone. As a result, corneal scarring from all causes is now the fourth greatest cause of global blindness.
Injuries
Eye injuries, most often occurring in people under 30, are the leading cause of monocular blindness (vision loss in one eye) throughout the United States. Injuries and cataracts affect the eye itself, while abnormalities such as optic nerve hypoplasia affect the nerve bundle that sends signals from the eye to the back of the brain, which can lead to decreased visual acuity.
Cortical blindness results from injuries to the occipital lobe of the brain that prevent the brain from correctly receiving or interpreting signals from the optic nerve. Symptoms of cortical blindness vary greatly across individuals and may be more severe in periods of exhaustion or stress. It is common for people with cortical blindness to have poorer vision later in the day.
Blinding has been used as an act of vengeance and torture in some instances, to deprive a person of a major sense by which they can navigate or interact within the world, act fully independently, and be aware of events surrounding them. An example from the classical realm is Oedipus, who gouges out his own eyes after realizing that he fulfilled the awful prophecy spoken of him. Having crushed the Bulgarians, the Byzantine Emperor Basil II blinded as many as 15,000 prisoners taken in the battle, before releasing them. Contemporary examples include the addition of methods such as acid throwing as a form of disfigurement.
Genetic defects
People with albinism often have vision loss to the extent that many are legally blind, though few of them actually cannot see. Leber congenital amaurosis can cause total blindness or severe sight loss from birth or early childhood. Retinitis pigmentosa is characterized by decreased peripheral vision and trouble seeing at night.
Advances in mapping of the human genome have identified other genetic causes of low vision or blindness. One such example is Bardet–Biedl syndrome.
Poisoning
Rarely, blindness is caused by the intake of certain chemicals. A well-known example is methanol, which is only mildly toxic and minimally intoxicating, and breaks down into the substances formaldehyde and formic acid which in turn can cause blindness, an array of other health complications, and death. When competing with ethanol for metabolism, ethanol is metabolized first, and the onset of toxicity is delayed. Methanol is commonly found in methylated spirits, denatured ethyl alcohol, to avoid paying taxes on selling ethanol intended for human consumption. Methylated spirits are sometimes used by alcoholics as a desperate and cheap substitute for regular ethanol alcoholic beverages.
Other
Amblyopia: is a category of vision loss or visual impairment that is caused by factors unrelated to refractive errors or coexisting ocular diseases. Amblyopia is the condition when a child's visual systems fail to mature normally because the child either has been born premature, measles, congenital rubella syndrome, vitamin A deficiency, or meningitis. If left untreated during childhood, amblyopia is currently incurable in adulthood because surgical treatment effectiveness changes as a child matures. Consequently, amblyopia is the world's leading cause of child monocular vision loss, which is the damage or loss of vision in one eye. In the best case scenario, which is very rare, properly treated amblyopia patients can regain 20/40 acuity.
Corneal opacification
Degenerative myopia
Diabetic retinopathy: is one of the manifestation microvascular complications of diabetes, which is characterized by blindness or reduced acuity. That is, diabetic retinopathy describes the retinal and vitreous hemorrhages or retinal capillary blockage caused by the increase of A1C, which a measurement of blood glucose or sugar level. In fact, as A1C increases, people tend to be at greater risk of developing diabetic retinopathy than developing other microvascular complications associated with diabetes (e.g. chronic hyperglycemia, diabetic neuropathy, and diabetic nephropathy). Despite the fact that only 8% of adults 40 years and older experience vision-threatening diabetic retinopathy (e.g. nonproliferative diabetic retinopathy or NPDR and proliferative diabetic retinopathy or PDR), this eye disease accounted for 17% of cases of blindness in 2002.
Retinitis pigmentosa
Retinopathy of prematurity: The most common cause of blindness in infants worldwide. In its most severe form, ROP causes retinal detachment, with attendant visual loss. Treatment is aimed mainly at prevention, via laser or Avastin therapy.
Stargardt's disease
Uveitis: is a group of 30 intraocular inflammatory diseases caused by infections, systemic diseases, organ-specific autoimmune processes, cancer or trauma. That is, uveitis refers to a complex category of ocular diseases that can cause blindness if either left untreated or improperly diagnosed. The current challenge of accurately diagnosing uveitis is that often the cause of a specific ocular inflammation is either unknown or multi-layered. Consequently, about 3–10% of those with uveitis in developed countries, and about 25% of those with uveitis in the developing countries, become blind from incorrect diagnosis and from ineffectual prescription of drugs, antibiotics or steroids. In addition, uveitis is a diverse category of eye diseases that are subdivided as granulomatous (or tumorous) or non-granulomatous anterior, intermediate, posterior or pan uveitis. In other words, uveitis diseases tend to be classified by their anatomic location in the eye (e.g. uveal tract, retina, or lens), as well as can create complication that can cause cataracts, glaucoma, retinal damage, age-related macular degeneration or diabetic retinopathy.
Xerophthalmia, often due to vitamin A deficiency, is estimated to affect 5 million children each year; 500,000 develop active corneal involvement, and half of these go blind.
Diagnosis
It is important that people be examined by someone specializing in low vision care prior to other rehabilitation training to rule out potential medical or surgical correction for the problem and to establish a careful baseline refraction and prescription of both normal and low vision glasses and optical aids. Only a doctor is qualified to evaluate visual functioning of a compromised visual system effectively. The American Medical Association provides an approach to evaluating visual loss as it affects an individual's ability to perform activities of daily living.
Screening adults who have no symptoms is of uncertain benefit.
Prevention
The World Health Organization estimates that 80% of visual loss is either preventable or curable with treatment. This includes cataracts, onchocerciasis, trachoma, glaucoma, diabetic retinopathy, uncorrected refractive errors, and some cases of childhood blindness. The Center for Disease Control and Prevention estimates that half of blindness in the United States is preventable.
Management
Mobility
Many people with serious visual impairments can travel independently, using a wide range of tools and techniques. Orientation and mobility specialists are professionals who are specifically trained to teach people with visual impairments how to travel safely, confidently, and independently in the home and the community. These professionals can also help blind people to practice travelling on specific routes which they may use often, such as the route from one's house to a convenience store. Becoming familiar with an environment or route can make it much easier for a blind person to navigate successfully.
Tools such as the white cane with a red tip – the international symbol of blindness – may also be used to improve mobility. A long cane is used to extend the user's range of touch sensation. It is usually swung in a low sweeping motion, across the intended path of travel, to detect obstacles. However, techniques for cane travel can vary depending on the user and/or the situation. Some visually impaired persons do not carry these kinds of canes, opting instead for the shorter, lighter identification (ID) cane. Still others require a support cane. The choice depends on the individual's vision, motivation, and other factors.
A small number of people employ guide dogs to assist in mobility. These dogs are trained to navigate around various obstacles, and to indicate when it becomes necessary to go up or down a step. However, the helpfulness of guide dogs is limited by the inability of dogs to understand complex directions. The human half of the guide dog team does the directing, based upon skills acquired through previous mobility training. In this sense, the handler might be likened to an aircraft's navigator, who must know how to get from one place to another, and the dog to the pilot, who gets them there safely.
GPS devices can also be used as a mobility aid. Such software can assist blind people with orientation and navigation, but it is not a replacement for traditional mobility tools such as white canes and guide dogs.
Some blind people are skilled at echolocating silent objects simply by producing mouth clicks and listening to the returning echoes. It has been shown that blind echolocation experts use what is normally the "visual" part of their brain to process the echoes.
Government actions are sometimes taken to make public places more accessible to blind people. Public transportation is freely available to blind people in many cities. Tactile paving and audible traffic signals can make it easier and safer for visually impaired pedestrians to cross streets. In addition to making rules about who can and cannot use a cane, some governments mandate the right-of-way be given to users of white canes or guide dogs.
Reading and magnification
Most visually impaired people who are not totally blind read print, either of a regular size or enlarged by magnification devices. Many also read large-print, which is easier for them to read without such devices. A variety of magnifying glasses, some handheld, and some on desktops, can make reading easier for them.
Others read braille (or the infrequently used Moon type), or rely on talking books and readers or reading machines, which convert printed text to speech or braille. They use computers with special hardware such as scanners and refreshable braille displays as well as software written specifically for the blind, such as optical character recognition applications and screen readers.
Some people access these materials through agencies for the blind, such as the National Library Service for the Blind and Physically Handicapped in the United States, the National Library for the Blind or the RNIB in the United Kingdom.
Closed-circuit televisions, equipment that enlarges and contrasts textual items, are a more high-tech alternative to traditional magnification devices.
There are also over 100 radio reading services throughout the world that provide people with vision impairments with readings from periodicals over the radio. The International Association of Audio Information Services provides links to all of these organizations.
Computers and mobile technology
Access technology such as screen readers, screen magnifiers and refreshable braille displays enable the blind to use mainstream computer applications and mobile phones. The availability of assistive technology is increasing, accompanied by concerted efforts to ensure the accessibility of information technology to all potential users, including the blind. Later versions of Microsoft Windows include an Accessibility Wizard & Magnifier for those with partial vision, and Microsoft Narrator, a simple screen reader. Linux distributions (as live CDs) for the blind include Vinux and Adriane Knoppix, the latter developed in part by Adriane Knopper who has a visual impairment. macOS and iOS also come with a built-in screen reader called VoiceOver, while Google TalkBack is built in to most Android devices.
The movement towards greater web accessibility is opening a far wider number of websites to adaptive technology, making the web a more inviting place for visually impaired surfers.
Experimental approaches in sensory substitution are beginning to provide access to arbitrary live views from a camera.
Modified visual output that includes large print and/or clear simple graphics can be of benefit to users with some residual vision.
Other aids and techniques
Blind people may use talking equipment such as thermometers, watches, clocks, scales, calculators, and compasses. They may also enlarge or mark dials on devices such as ovens and thermostats to make them usable. Other techniques used by blind people to assist them in daily activities include:
Adaptations of coins and banknotes so that the value can be determined by touch. For example:
In some currencies, such as the euro, the pound sterling and the Indian rupee, the size of a note increases with its value.
On US coins, pennies and dimes, and nickels and quarters are similar in size. The larger denominations (dimes and quarters) have ridges along the sides (historically used to prevent the "shaving" of precious metals from the coins), which can now be used for identification.
Some currencies' banknotes have a tactile feature to indicate denomination. For example, the Canadian currency tactile feature is a system of raised dots in one corner, based on braille cells but not standard braille.
It is also possible to fold notes in different ways to assist recognition.
Labeling and tagging clothing and other personal items
Placing different types of food at different positions on a dinner plate
Marking controls of household appliances
Most people, once they have been visually impaired for long enough, devise their own adaptive strategies in all areas of personal and professional management.
For the blind, there are books in braille, audio-books, and text-to-speech computer programs, machines and e-book readers. Low vision people can make use of these tools as well as large-print reading materials and e-book readers that provide large font sizes.
Computers are important tools of integration for the visually impaired person. They allow, using standard or specific programs, screen magnification and conversion of text into sound or touch (braille line), and are useful for all levels of visual impairment. OCR scanners can, in conjunction with text-to-speech software, read the contents of books and documents aloud via computer. Vendors also build closed-circuit televisions that electronically magnify paper, and even change its contrast and color, for visually impaired users. For more information, consult assistive technology.
In adults with low vision there is no conclusive evidence supporting one form of reading aid over another. In several studies stand-mounted devices allowed faster reading than hand-held or portable optical aids. While electronic aids may allow faster reading for individuals with low vision, portability, ease of use, and affordability must be considered for people.
Children with low vision sometimes have reading delays, but do benefit from phonics-based beginning reading instruction methods. Engaging phonics instruction is multisensory, highly motivating, and hands-on. Typically students are first taught the most frequent sounds of the alphabet letters, especially the so-called short vowel sounds, then taught to blend sounds together with three-letter consonant-vowel-consonant words such as cat, red, sit, hot, sun. Hands-on (or kinesthetically appealing) VERY enlarged print materials such as those found in "The Big Collection of Phonics Flipbooks" by Lynn Gordon (Scholastic, 2010) are helpful for teaching word families and blending skills to beginning readers with low vision. Beginning reading instructional materials should focus primarily on the lower-case letters, not the capital letters (even though they are larger) because reading text requires familiarity (mostly) with lower-case letters. Phonics-based beginning reading should also be supplemented with phonemic awareness lessons, writing opportunities, and many read-alouds (literature read to children daily) to stimulate motivation, vocabulary development, concept development, and comprehension skill development. Many children with low vision can be successfully included in regular education environments. Parents may need to be vigilant to ensure that the school provides the teacher and students with appropriate low vision resources, for example technology in the classroom, classroom aide time, modified educational materials, and consultation assistance with low vision experts.
Epidemiology
The WHO estimates that in 2012 there were 285 million visually impaired people in the world, of which 246 million had low vision and 39 million were blind.
Of those who are blind 90% live in the developing world. Worldwide for each blind person, an average of 3.4 people have low vision, with country and regional variation ranging from 2.4 to 5.5.
By age: Visual impairment is unequally distributed across age groups. More than 82% of all people who are blind are 50 years of age and older, although they represent only 19% of the world's population. Due to the expected number of years lived in blindness (blind years), childhood blindness remains a significant problem, with an estimated 1.4 million blind children below age 15.
By gender: Available studies consistently indicate that in every region of the world, and at all ages, females have a significantly higher risk of being visually impaired than males.
By geography: Visual impairment is not distributed uniformly throughout the world. More than 90% of the world's visually impaired live in developing countries.
Since the estimates of the 1990s, new data based on the 2002 global population show a reduction in the number of people who are blind or visually impaired, and those who are blind from the effects of infectious diseases, but an increase in the number of people who are blind from conditions related to longer life spans.
In 1987, it was estimated that 598,000 people in the United States met the legal definition of blindness. Of this number, 58% were over the age of 65. In 1994–1995, 1.3 million Americans reported legal blindness.
Society and culture
Legal definition
To determine which people qualify for special assistance because of their visual disabilities, various governments have specific definitions for legal blindness. In North America and most of Europe, legal blindness is defined as visual acuity (vision) of 20/200 (6/60) or less in the better eye with best correction possible. This means that a legally blind individual would have to stand from an object to see itwith corrective lenseswith the same degree of clarity as a normally sighted person could from . In many areas, people with average acuity who nonetheless have a visual field of less than 20 degrees (the norm being 180 degrees) are also classified as being legally blind.
Approximately fifteen percent of those deemed legally blind, by any measure, have no light or form perception. The rest have some vision, from light perception alone to relatively good acuity. Low vision is sometimes used to describe visual acuities from 20/70 to 20/200.
Literature and art
Antiquity
The Moche people of ancient Peru depicted the blind in their ceramics.
In Greek myth, Tiresias was a prophet famous for his clairvoyance. According to one myth, he was blinded by the gods as punishment for revealing their secrets, while another holds that he was blinded as punishment after he saw Athena naked while she was bathing. In the Odyssey, the one-eyed Cyclops Polyphemus captures Odysseus, who blinds Polyphemus to escape. In Norse mythology, Loki tricks the blind god Höðr into killing his brother Baldr, the god of happiness.
The New Testament contains numerous instances of Jesus performing miracles to heal the blind. According to the Gospels, Jesus healed the two blind men of Galilee, the blind man of Bethsaida, the blind man of Jericho and the man who was born blind.
The parable of the blind men and an elephant has crossed between many religious traditions and is part of Jain, Buddhist, Sufi and Hindu lore. In various versions of the tale, a group of blind men (or men in the dark) touch an elephant to learn what it is like. Each one feels a different part, but only one part, such as the side or the tusk. They then compare notes and learn that they are in complete disagreement.
"Three Blind Mice" is a medieval English nursery rhyme about three blind mice whose tails are cut off after chasing the farmer's wife. The work is explicitly incongruous, ending with the comment Did you ever see such a sight in your life, As three blind mice?
Modern times
Poet John Milton, who went blind in mid-life, composed "On His Blindness", a sonnet about coping with blindness. The work posits that [those] who best Bear [God]'s mild yoke, they serve him best.
The Dutch painter and engraver Rembrandt often depicted scenes from the apocryphal Book of Tobit, which tells the story of a blind patriarch who is healed by his son, Tobias, with the help of the archangel Raphael.
Slaver-turned-abolitionist John Newton composed the hymn "Amazing Grace" about a wretch who "once was lost, but now am found, Was blind, but now I see." Blindness, in this sense, is used both metaphorically (to refer to someone who was ignorant but later became knowledgeable) and literally, as a reference to those healed in the Bible. In the later years of his life, Newton himself would go blind.
H. G. Wells' story "The Country of the Blind" explores what would happen if a sighted man found himself trapped in a country of blind people to emphasise society's attitude to blind people by turning the situation on its head.
José Saramago's novel Blindness describes a sudden mass blindness epidemic, and the main story follows the first group to lose their vision and how they have to adapt to their new disability amidst a decaying society.
Bob Dylan's anti-war song "Blowin' in the Wind" twice alludes to metaphorical blindness: How many times can a man turn his head // and pretend that he just doesn't see... How many times must a man look up // Before he can see the sky?
Contemporary fiction contains numerous well-known blind characters. Some of these characters can see by means of devices, such as the Marvel Comics superhero Daredevil, who can see via his super-human hearing acuity, or Star Trek Geordi La Forge, who can see with the aid of a VISOR, a fictional device that transmits optical signals to his brain.
Blind culture
People who are not blind often imagine that people who are blind share a cultural identity in the way that other minority groups with shared experiences have a distinct culture. Various blind commentators have responded to this perception by explaining that more commonly, blind people integrate with the broader community and culture, and often do not identify blindness as a defining part of their culture.
People who are blind share the common cultural experience of the many misconceptions sighted people have about living with blindness.
Sports
Blind and partially sighted people participate in sports, such as swimming, snow skiing and athletics. Some sports have been invented or adapted for the blind, such as goalball, association football, cricket, golf, tennis, bowling, and beep baseball. The worldwide authority on sports for the blind is the International Blind Sports Federation. People with vision impairments have participated in the Paralympic Games since the 1976 Toronto summer Paralympics.
Metaphorical uses
The word "blind" (adjective and verb) is often used to signify a lack of knowledge of something. For example, a blind date is a date in which the people involved have not previously met; a blind experiment is one in which information is kept from either the experimenter or the participant to mitigate the placebo effect or observer bias. The expression "blind leading the blind" refers to incapable people leading other incapable people. Being blind to something means not understanding or being aware of it. A "blind spot" is an area where someone cannot see: for example, where a car driver cannot see because parts of his car's bodywork are in the way; metaphorically, a topic on which an individual is unaware of their own biases, and therefore of the resulting distortions of their own judgements (see Bias blind spot).
Research
A 2008 study tested the effect of using gene therapy to help restore the sight of patients with a rare form of inherited blindness, known as Leber's congenital amaurosis or LCA. Leber's Congenital Amaurosis damages the light receptors in the retina and usually begins affecting sight in early childhood, with worsening vision until complete blindness around the age of 30. The study used a common cold virus to deliver a normal version of the gene called RPE65 directly into the eyes of affected patients. All three patients, aged 19, 22 and 25, responded well to the treatment and reported improved vision following the procedure.
Two experimental treatments for retinal problems include a cybernetic replacement and transplant of fetal retinal cells.
There is no high-quality evidence on the effect of assistive technologies on educational outcomes and quality of life in children with low vision , nor is there evidence on magnifying reading aids in children. Low-vision rehabilitation does not appear to have an important impact on health-related quality of life, though some low-vision rehabilitation interventions, particularly psychological therapies and methods of enhancing vision, may improve vision-related quality of life in people with sight loss.
Other animals
Statements that certain species of mammals are "born blind" refers to them being born with their eyes closed and their eyelids fused together; the eyes open later. One example is the rabbit. In humans, the eyelids are fused for a while before birth, but open again before the normal birth time; however, very premature babies are sometimes born with their eyes fused shut, and opening later. Other animals, such as the blind mole rat, are truly blind and rely on other senses.
The theme of blind animals has been a powerful one in literature. Peter Shaffer's Tony Award-winning play, Equus, tells the story of a boy who blinds six horses. Theodore Taylor's classic young adult novel, The Trouble With Tuck, is about a teenage girl, Helen, who trains her blind dog to follow and trust a seeing-eye dog.
| Biology and health sciences | Disability | null |
2158876 | https://en.wikipedia.org/wiki/Port%20of%20Hamburg | Port of Hamburg | The Port of Hamburg (, ) is a seaport on the river Elbe in Hamburg, Germany, from its mouth on the North Sea.
Known as Germany's "Gateway to the World" (), it is the country's largest seaport by volume. In terms of TEU throughput, Hamburg is the third-busiest port in Europe (after Rotterdam and Antwerp) and 15th-largest worldwide. In 2014, 9.73 million TEUs (20-foot standard container equivalents) were handled in Hamburg.
The port covers an area of (64.80 km2 usable), of which 43.31 km2 (34.12 km2) are land areas. The branching Elbe creates an ideal place for a port complex with warehousing and transshipment facilities. The extensive free port was established when Hamburg joined the German Customs Union. It enabled duty-free storing of imported goods and also importing of materials which were processed, re-packaged, used in manufacturing and then re-exported without incurring customs duties. The free port was abandoned in 2013.
History
The port is almost as old as the history of Hamburg itself. Founded on 7 May 1189 by Frederick I at a strategic location near the mouth of the Elbe, it has been Central Europe's main port for centuries and enabled Hamburg to develop early into a leading city of trade with a rich and proud bourgeoisie.
During the age of the Hanseatic League from the 13th to 16th century, Hamburg was considered second only to the port and city of Lübeck in terms of its position as a central trading node for sea-borne trade. With discovery of the Americas and the emerging transatlantic trade, Hamburg exceeded all other German ports. During the second half of the 19th century, Hamburg became Central Europe's main hub for transatlantic passenger and freight travel, and from 1871 onward it was Germany's principal port of trade. In her time the Hamburg America Line was the largest shipping company in the world. Since 1888, the HADAG runs a scheduled ferry service across various parts of the port and the Elbe. The Free Port (Freihafen), established on 15 October 1888, enabled traders to ship and store goods without going through customs and further enhanced Hamburg's position in sea trade with neighbouring countries. It was permanently closed on 1 January 2013. The Moldauhafen has a similar arrangement, though related to the Czech Republic exclusively.
The Speicherstadt, one of Hamburg's architectural icons today, is a large wharf area of 350,000 m2 floor area on the northern shore of the river, built in the 1880s as part of the free port and to cope with the growing quantity of goods stored in the port.
Hamburg shipyards lost fleets twice after World War I and World War II. Moreover, during World War II, Hamburg harbour was the hub destination of the Hamburg America Line, that assured the Nazi Party a connection to the United States for the import of oil and steel, and the export of manufactured goods from Germany thanks to container ships. The shipping line Hamburg-Amerikanische Packetfahrt-Aktien-Gesellschaft (HAPAG) gave the name to the so-called shipping company based in Hamburg which used to run the trades of goods on this route. In 1970, along with Norddeutscher Lloyd, the present-day active company Hapag-Lloyd was founded.
During the partition of Germany between 1945 and 1990, the Port of Hamburg lost much of its hinterland and consequently many of its trading connections. However, since German reunification, the fall of the Iron Curtain and European enlargement, Hamburg has made substantial ground as one of Europe's prime logistics centres and as one of the world's largest and busiest sea ports.
In 2022, the German government let the Chinese state-owned COSCO Shipping take a stake in ownership of the port.
Access
Deepening of the river Elbe for large vessels is controversial for ecological reasons. In part due to cooperation with Lower Saxony and Bremen to build a new container port (JadeWeserPort) in the deep waters of Jadebusen in Wilhelmshaven, Hamburg withdrew from this plan after a change of government in 2001.
Hamburg Port Authority
The port is administered by the Hamburg Port Authority. The Hamburg Port Authority is described as having adopted an innovative approach. In November 2016 the Hamburg Port Authority ordered a modern fireboat budgeted at 16 million euros.
Terminals
Cruise
Hamburg is a major cruise destination and one of Europe's largest ports of call for cruise passengers traveling the Atlantic, or the Norwegian and Baltic Seas. The port is also a major location for shipbuilder and shipyards, designing, building and reconditioning yachts and cruise liners. Hamburg has three passenger terminals for cruise ships: Hamburg Cruise Center HafenCity, the Hamburg Cruise Center Altona and the Hamburg Cruise Center Steinwerder, all three capable of processing the world's largest cruise ships.
Culture
Hamburg's harbour is also one of the city's major attractions, both as a vital, industrial, and logistical centre, and as a backdrop for modern culture and harbour history. These include several museum ships, musical theatres, bars, restaurants, and hotels – and even a floating church.
The annual celebration of the port's birthday (Hafengeburtstag), during the first weekend of May, is one of Hamburg's biggest public events. National and international visitors come to experience the festivities. Tugboats perform "ballets", old galleons and new cruise ships are open for tours, and fireworks explode at night.
Tour guides on boat tours in the port are called "he lüchts" (Low German for "he is lying"), after an often used call of dock workers when they overheard the stories told to tourists.
| Technology | Specific piers and ports | null |
2159778 | https://en.wikipedia.org/wiki/Sexual%20and%20reproductive%20health | Sexual and reproductive health | Sexual and reproductive health (SRH) is a field of research, health care, and social activism that explores the health of an individual's reproductive system and sexual well-being during all stages of their life. Sexual and reproductive health is more commonly defined as sexual and reproductive health and rights, to encompass individual agency to make choices about their sexual and reproductive lives.
The term can also be further defined more broadly within the framework of the World Health Organization's (WHO) definition of health―as "a state of complete physical, mental and social well-being, and not merely the absence of disease or infirmity"―. WHO has a working definition of sexual health (2006) as '“…a state of physical, emotional, mental and social well-being in relation to sexuality; it is not merely the absence of disease, dysfunction or infirmity. Sexual health requires a positive and respectful approach to sexuality and sexual relationships, as well as the possibility of having pleasurable and safe sexual experiences, free of coercion, discrimination and violence. For sexual health to be attained and maintained, the sexual rights of all persons must be respected, protected and fulfilled.” This includes sexual wellbeing, encompassing the ability of an individual to have responsible, satisfying and safe sex and the freedom to decide if, when and how often to do so. UN agencies in particular define sexual and reproductive health as including both physical and psychological well-being vis-à-vis sexuality. Furthermore, the importance of ensuring sexual lives are pleasurable and satisfying, and not only focused on negative consequences of sex has been emphasized by many agencies such as the World Association of Sexual Health as well as considering the positive impacts on health and well-being of safe and satisfying relationships. A further interpretation includes access to sex education, access to safe, effective, affordable and acceptable methods of birth control, as well as access to appropriate health care services, as the ability of women to go safely through pregnancy and childbirth could provide couples with the best chance of having a healthy infant.
The critical Guttmacher- Lancet Commission on Sexual and reproductive health and rights states state Sexual and reproductive health and rights (SRHR) are essential for sustainable development because of their links to gender equality and women’s wellbeing, their impact on maternal, newborn, child, and adolescent health, and their roles in shaping future economic development and environmental sustainability. Yet progress towards fulfilling SRHR for all has been stymied because of weak political commitment, inadequate resources, persistent discrimination against women and girls, and an unwillingness to address issues related to sexuality openly and comprehensively. As a result, almost all of the 4·3 billion people of reproductive age worldwide will have inadequate sexual and reproductive health services over the course of their lives'.
Individuals face inequalities in reproductive health services. Inequalities vary based on socioeconomic status, education level, age, ethnicity, religion, and resources available in their environment. Low income individuals may lack access to appropriate health services and/or knowledge of how to maintain reproductive health. Additionally, many approaches involving women, families, and local communities as active stakeholders in interventions and strategies to improve reproductive health.
Overview
The WHO assessed in 2008 that "Reproductive and sexual ill-health accounts for 20% of the global burden of ill-health for women, and 14% for men." Reproductive health is a part of sexual and reproductive health and rights. According to the United Nations Population Fund (UNFPA), unmet needs for sexual and reproductive health deprive women of the right to make "crucial choices about their own bodies and futures", affecting family welfare. Women bear and usually nurture children, so their reproductive health is inseparable from gender equality. Denial of such rights also worsens poverty.
Adolescent health
Adolescent health creates a major global burden and has a great deal of additional and diverse complications compared to adult reproductive health such as early pregnancy and parenting issues, difficulties accessing contraception and safe abortions, lack of healthcare access, and high rates of HIV, sexually transmitted infections and mental health issues. Each of those can be affected by outside political, economic and socio-cultural influences. For most adolescent females, they have yet to complete their body growth trajectories, therefore adding a pregnancy exposes them to a predisposition to complications. These complications range from anemia, malaria, HIV and other STIs, postpartum bleeding and other postpartum complications, mental health disorders such as depression and suicidal thoughts or attempts. In 2016, adolescent birth rates between the ages of 15-19 was 45 per 1000. In 2014, 1 in 3 experienced sexual violence, and there more than 1.2 million deaths. The top three leading causes of death in females between the ages of 15-19 are maternal conditions 10.1%, self-harm 9.6%, and road conditions 6.1%.
The causes of teenage pregnancy are vast and diverse. In developing countries, young women are pressured to marry for different reasons. One reason is to bear children to help with work, another on a dowry system to increase the families income, another is due to prearranged marriages. These reasons tie back to the financial needs of girls' families, cultural norms, religious beliefs, and external conflicts.
Adolescent pregnancy, especially in developing countries, carries increased health risks, and contributes to maintaining the cycle of poverty. The availability and type of sex education for teenagers varies in different parts of the world. Teens that are self-identified as non-heterosexual may develop additional problems if they live in places where homosexual activity is socially disapproved or even illegal; in extreme cases, there can be depression, social isolation, and even suicide among LGBT youth.
Maternal health
95% of maternal deaths occur in low income contexts and countries, and in 25 years, the maternal mortality globally dropped to 44%. Statistically, a woman's chance of survival during childbirth is closely tied to her social-economic status, access to healthcare, where she lives geographically, and cultural norms. To compare, a woman dies of complications from childbirth every minute in developing countries versus a total of 1% of total maternal mortality deaths in developed countries. Women in developing countries have little access to family planning services, different cultural practices, lack of information, birthing attendants, prenatal care, birth control, postnatal care, lack of access to health care, and are typically in poverty. In 2015, those in low-income countries had access to antenatal care visits averaged 40% and were preventable. All these reasons led to an increase in the maternal mortality ratio (MMR).
One of the international Sustainable Development Goals developed by United Nations is to improve maternal health by a targeted 70 deaths per 100,000 live births by 2030. Most models of maternal health encompass family planning, preconception, prenatal, and postnatal care. All care after childbirth recovery is typically excluded, which includes pre-menopause and aging into old age. During childbirth, women typically die from severe bleeding, infections, high blood pressure during pregnancy, delivery complications, or an unsafe abortion. Other reasons can be regional such as complications related to diseases such as malaria and AIDS during pregnancy. The younger the woman is when she gives birth, the more at risk she and her baby are for complications and possible mortality.
There is a significant relationship between the quality of maternal services made available and the greater financial standings of a country. Sub-Saharan Africa and South Asia exemplify this as these regions are significantly deprived of medical staff and affordable health
opportunities. Most countries provide for their health services through a combination of funding from government tax revenue and local households. Poorer nations or regions with extremely concentrated wealth can leave citizens on the margins uncared for or overlooked.
However, the lack of proper leadership can result in a nation's public sectors being mishandled or poorly performing despite said nation's resources and standing. In addition, poorer nations funding their medical services through taxes places a greater financial burden on the public and effectively the mothers themselves.
Responsibility and accountability on the part of mental health sectors are strongly emphasized as to what will remedy the poor quality of maternal health globally. The impact of different maternal health interventions across the globe stagger variously and are vastly uneven. This is the result of a lack of political and financial commitment to the issue as most safe motherhood programs internationally have to compete for significant funding. Some resolve that if global survival initiatives were promoted and properly funded it would prove to be mutually beneficial for the international community. Investing in maternal health would ultimately advance several issues such as gender inequality, poverty, and general global health standards. As it currently stands, pregnant women are subjugated to high financial costs throughout the duration of their term internationally which is highly taxing and strenuous.
In addition, if either parent has a genetic disease, there is risk of these being passed on to the children. Birth control or technical solutions (assisted reproductive technology) can be an option then.
LGBT+ sexual and reproductive health
The sexual and reproductive health of LGBT+ people face challenges through issues like the ongoing HIV pandemic, binary organization of "men" and "women"'s reproductive health, alongside stigma and repression that limit LGBT+ people from accessing the healthcare they need. Sexual health is a state of physical, emotional, mental, and social well-being in relation to sexuality. It is important to not only consider the sexual/physical health of an individual but also the cultural and contextual factors that influence an individual's well-being. The lack of competent providers and stigma attached to homosexuality have a great impact on the sexual health of the LGBTQ+ population. The LGBTQ+ population faces a number of obstacles in terms of sexual and reproductive health. The different stigmas and biases that come with these barriers make receiving proper care difficult. Some of these stigmas that follow those in the LGBTQ+ population in terms of their sexual and reproductive health are associating certain diseases, and other illnesses with this community. This leaves those in the LGBTQ+ population in a position that makes them vulnerable, as well as victims of a number of health disparities. The overall health of those in the LGBTQ+ population is determinant on sexual and reproductive health as these all make up the health of these individuals. Those in the LGBTQ+ community face also face discrimination from providers and insurance companies, on top of all of the other barriers and limits on access to care that they endure. All of these factors have led to those in the LGBTQ+ population having worse health outcomes.
Contraception
Access to reproductive health services is very poor in many countries. Women are often unable to access maternal health services due to a lack of knowledge about the existence of such services or lack of freedom of movement. Some women are subjected to forced pregnancy and banned from leaving the home. In many countries, women are not allowed to leave home without a male relative or husband, and therefore their ability to access medical services is limited. Therefore, increasing women's autonomy is needed in order to improve reproductive health, however, doing so may require a cultural shift. According to the WHO, "All women need access to antenatal care in pregnancy, skilled care during childbirth, and care and support in the weeks after childbirth".
The fact that the law allows certain reproductive health services, does not necessarily ensure that such services are actually in use by the people. The availability of contraception, sterilization, and abortion is dependent on laws, as well as social, cultural, and religious norms. Some countries have liberal laws regarding these issues, but in practice, it is very difficult to access such services due to doctors, pharmacists, and other social and medical workers being conscientious objectors.
In developing regions of the world, there are about 214 million women who want to avoid pregnancy but are unable to use safe and effective family planning methods. When taken correctly, the combined oral contraceptive pill is over 99% effective at preventing pregnancy. However, it does not protect from sexually transmitted infections (STIs). Some methods, such as using condoms, achieve both protection from STIs and unwanted pregnancies. There are also natural family planning methods, which may be preferred by religious people, but some very conservative religious groups, such as the Quiverfull movement, oppose these methods too because they advocate the maximization of procreation. One of the oldest ways to reduce unwanted pregnancy is coitus interruptus - still widely used in the developing world.
There are many types of contraceptives. One type of contraceptive includes barrier methods. One barrier method includes condoms for males and females. Both types stop sperm from entering the woman's uterus, thereby preventing pregnancy from occurring. Another type of contraception is the birth control pill, which stops ovulation from occurring by combining the chemicals progestin and estrogen. Many women use this method of contraception, however, they discontinue using it equally as much as they use it. One reason for this is because of the side effects that may occur from using the pill, and because some health care providers do not take women's concerns about negative side effects seriously. The use of the birth control pill is common in western countries, and two forms of combined oral contraceptives are on the World Health Organization's List of Essential Medicines, the most important medications needed in a basic health system.
There are many objections to the use of birth control, both historically and in the present day. One argument against birth control usage states that there is no need for birth control, to begin with. This argument was levied in 1968 when Richard Nixon was elected president, and the argument stated that since birth rates were at their lowest point since World War II ended, birth control was not necessary. Demographic planning arguments were also the basis of the population policy of Nicolae Ceaușescu in communist Romania, who adopted a very aggressive natalist policy which included outlawing abortion and contraception, routine pregnancy tests for women, taxes on childlessness, and legal discrimination against childless people. Such policies consider that coercion is an acceptable means of reaching demographic targets. Religious objections are based on the view that premarital sex should not happen, while married couples should have as many children as possible. As such, the Catholic Church encourages premarital abstinence from sex. This argument was written out in Humanae Vitae, a papal encyclical released in 1968. The Catholic Church bases its argument against birth control pills on the basis that birth control pills undermine the natural law of God. The Catholic Church also argues against birth control on the basis of family size, with Cardinal Mercier of Belgium arguing, "...the duties of conscience are above worldly considerations, and besides, it is the large families who are the best" (Reiterman, 216). Another argument states that women should use natural methods of contraception in place of artificial ones, such as having sexual intercourse when one is infertile.
Support for contraception is based on views such as reproductive rights, women's rights, and the necessity to prevent child abandonment and child poverty. The World Health Organization states that "By preventing unintended pregnancy, family planning /contraception prevents deaths of mothers and children".
Sexually transmitted infection
A sexually transmitted infection (STI) --previously known as a sexually transmitted disease (STD) or venereal disease (VD)-- is an infection that has a significant likelihood of transmission between humans by means of sexual activity. The CDC analyses the eight most common STIs: chlamydia, gonorrhea, hepatitis B virus (HBV), herpes simplex virus type 2 (HSV-2), human immunodeficiency virus (HIV), human papillomavirus (HPV), syphilis, and trichomoniasis.
There are 1 million new infections a day and more than 20 million new cases within the United States. In 2020, WHO estimated 374 million new infections with 1 of 4 STIs: chlamydia (129 million), gonorrhoea (82 million), syphilis (7.1 million) and trichomoniasis (156 million). More than 490 million people were estimated to be living with genital herpes in 2016, and an estimated 300 million women have an HPV infection, the primary cause of cervical cancer and anal cancer among men who have sex with men.
Numbers of such high magnitude weigh a heavy burden on the local and global economy. A study conducted at Oxford University in 2015 concluded that despite giving participants early antiviral medications (ART), they still cost an estimated $256 billion over 2 decades. HIV testing done at modest rates could reduce HIV infections by 21%, HIV retention by 54%, and HIV mortality rates by 64%, with a cost-effectiveness ratio of $45,300 per quality-adjusted life year. However, the study concluded that the United States has led to an excess in infections, treatment costs, and deaths, even when interventions do not improve overall survival rates.
There is a profound reduction in STI rates once those who are sexually active are educated about transmissions, condom promotion, interventions targeted at key and vulnerable populations through comprehensive sex education courses or programs. Recent evidence shows that acknowledging the role pleasure takes in people's sexual lives and integrating this in sexual health services and education has a significant impact on increasing condom use and improved sexual health outcomes. South Africa's policy addresses the needs of women at risk for HIV and who are HIV positive as well as their partners and children. The policy also promotes screening activities related to sexual health such as HIV counseling and testing as well as testing for other STIs, tuberculosis, cervical cancer, and breast cancer.
The CDC stated that the rate of sexually transmitted infections is higher among minorities compared to white people. These minorities are currently being affected by different factors including health literacy, socioeconomic status, access to health services, and fear of discrimination by health providers. The rates of infection are five to eight times higher in the Black community compared to non-Hispanic White people.
Young African American women are at a higher risk for STIs, including HIV. A recent study published outside of Atlanta, Georgia collected data (demographic, psychological, and behavioral measures) with a vaginal swab to confirm the presence of STIs. They found a profound difference that those women who had graduated from college were far less likely to have STIs, potentially be benefiting from a reduction in vulnerability to acquiring STIs/HIV as they gain in education status and potentially move up in demographic areas and/or status.
Abortion
Globally, an estimated 25 million unsafe abortions occur each year. The vast majority of such unsafe abortions occur in developing countries in Africa, Asia and Latin America.
The abortion debate is the ongoing controversy surrounding the moral, legal, and religious status of induced abortion. The sides involved in the debate are the self-described "pro-choice" and "pro-life" movements. "Pro-choice" emphasizes the right of women to decide whether to terminate a pregnancy. "Pro-life" emphasizes the right of the embryo or fetus to gestate to term and be born. Both terms are considered loaded in mainstream media, where terms such as "abortion rights" or "anti-abortion" are generally preferred. Each movement has, with varying results, sought to influence public opinion and to attain legal support for its position, with small numbers of radical activists using violence, such as murder and arson.
Articles from the World Health Organization call legal abortion a fundamental right of women regardless of where they live, and argue that unsafe abortion is a silent pandemic. In 2005, it was estimated that 19-20 million abortions had complications, some complications are permanent, while another estimated 68,000 women died from unsafe abortions. Having access to safe abortion can have positive impacts on women's health and life, and vice versa. "Legislation of abortion on request is necessary but an insufficient step towards improving women's health. In some countries where it abortion is legal and has been for decades, there has been no improvement in access to adequate services making abortion unsafe due to lack of healthcare services. It is hard to get an abortion due to legal and policy barriers, social and cultural barriers (gender discrimination, poverty, religious restrictions, lack of support), health system barriers (lack of facilities or trained personnel). However, safe abortions with trained personnel, good social support, and access to facilities, can improve maternal health and increase reproductive health later in life.
The Maputo Protocol, which was adopted by the African Union in the form of a protocol to the African Charter on Human and Peoples' Rights, states at Article 14 (Health and Reproductive Rights) that: "(2). States Parties shall take all appropriate measures to: [...] c) protect the reproductive rights of women by authorising medical abortion in cases of sexual assault, rape, incest, and where the continued pregnancy endangers the mental and physical health of the mother or the life of the mother or the foetus." The Maputo Protocol is the first international treaty to recognize abortion, under certain conditions, as a woman's human right.
The General comment No. 36 (2018) on article 6 of the International Covenant on Civil and Political Rights, on the right to life, adopted by the Human Rights Committee in 2018, defines, for the first time ever, a human right to abortion - in certain circumstances (however these UN general comments are considered soft law, and, as such, not legally binding).
When negotiating the Cairo Programme of Action at the 1994 International Conference on Population and Development (ICPD), the issue was so contentious that delegates eventually decided to omit any recommendation to legalize abortion, instead advising governments to provide proper post-abortion care and to invest in programs that will decrease the number of unwanted pregnancies.
The Committee on the Elimination of Discrimination against Women considers the criminalization of abortion a "violations of women's sexual and reproductive health and rights" and a form of "gender based violence"; paragraph 18 of its General recommendation No. 35 on gender based violence against women, updating general recommendation No. 19 states that: "Violations of women's sexual and reproductive health and rights, such as forced sterilizations, forced abortion, forced pregnancy, criminalisation of abortion, denial or delay of safe abortion and post abortion care, forced continuation of pregnancy, abuse and mistreatment of women and girls seeking sexual and reproductive health information, goods and services, are forms of gender based violence that, depending on the circumstances, may amount to torture or cruel, inhuman or degrading treatment." The same General Recommendation also urges countries at paragraph 31 to [...] In particular, repeal:
a) Provisions that allow, tolerate or condone forms of gender based violence against women, including [...] legislation that criminalises abortion".
In 2008, the Parliamentary Assembly of the Council of Europe, a group comprising members from 47 European countries, has adopted a resolution calling for the decriminalization of abortion within reasonable gestational limits and guaranteed access to safe abortion procedures. The nonbinding resolution was passed on April 16 by a vote of 102 to 69.
Accesses to abortion is not only a question of legality but also an issue of overcoming de facto barriers, such as conscientious objections from medical staff, high prices, lack of knowledge about the law, lack of access to medical care (especially in rural areas). The de facto inability of women to access abortion even in countries where it is legal is highly controversial because it results in a situation where women have rights only on paper, not in practice; the UN in its 2017 resolution on Intensification of efforts to prevent and eliminate all forms of violence against women and girls: domestic violence urged states to guarantee access to "safe abortion where such services are permitted by national law".
There are two primary arguments for maintaining legalized abortion today in the U.S. The first is recognizing the full citizenship of women. The Roe v. Wade'' court case on abortion compared the citizenship of women and fetuses Because the Constitution defines born people as citizens, Justice Harry Blackmun ruled that fetuses were not citizens. The citizenship of women is emphasized because fetuses are not individual entities that can exist without the woman. Another reason why the full citizenship of women is defined by advocates for abortion is that it recognizes the right of women to manage their own bodies. Fertility affects women's bodies. The argument for abortion prevents others from making decisions that alter a woman's body. Pro-choice advocates also attempt to confirm that state-mandated education or other outside biases do not attempt to influence these decisions. Feminists argue that women throughout history have had to justify their citizenship politically and socially. The right to manage one's own body is a matter of health, safety, and respect. The citizenship of women and the right to manage their own bodies is a societal confirmation that feminists highlight as a pro-choice justification.
The second primary argument to uphold legalized abortion and creating better access to it is the necessity of abortion and the health and safety of pregnant women. There are two events that largely changed the course of public opinion about abortion in the U.S. The first is Sherry Finkbine, who was denied access to an abortion by the board of obstetrician-gynecologists at her local hospital. Although she was privileged enough to afford the trip, Finkbine was forced to travel to Sweden for an abortion to avoid caring for a damaged fetus in addition to four children. The other event that changed public opinion was the outbreak of rubella in the 1950s and 60s. Because rubella disrupted the growth of fetuses and caused deformities during pregnancy, the California Therapeutic Abortion Act was signed in 1967, permitting doctors to legally abort pregnancies that pose a risk to a pregnant woman's physical or mental health. These two events are commonly used to show how the health and safety of pregnant women are contingent upon abortions as well as the ability to give birth to and adequately take care of a child. Another argument in favor of legalized abortion to service necessity are the reasons why an abortion might be necessary. Nearly half of all pregnancies in the United States are unintended, and over half of all unintended pregnancies in the United States are met with abortion. Unintended pregnancy can lead to serious harm to women and children for reasons such as not being able to afford to raise a baby, inaccessibility to time off of work, difficulties facing single motherhood, difficult socio-economic conditions for women. Unintended pregnancies also have a greater potential for putting women of color at risk due to systematically produced environmental hazards from proximity to pollution, access to livable income, and affordable healthy food. These factors as threats to the health and safety of pregnant women run parallel to data that shows the number of abortions in the United States did not decline while laws restricting legal access to abortion were implemented.
At a global level, the region with the strictest abortion laws is considered to be Latin America (see Reproductive rights in Latin America), a region strongly influenced by the Catholic Church in Latin America.
Female genital mutilation
Female genital mutilation (FGM), also known as female genital circumcision or cutting, is the traditional, non-medical practice of altering or injuring the female reproductive organs, often by removing all or parts of the external genitalia. It is mostly practiced in 30 countries in Africa, the Middle East, and Asia, and affects over 200 million women and girls worldwide. More severe forms of FGM are highly concentrated in Djibouti, Eritrea, Ethiopia, Somalia, and Sudan.
The WHO categorizes FGM into four types:
Type I (Clitoridectomy) is the removal of all or part of the clitoris. This may or may not include removing the prepuce along with the clitoral glans.
Type II (Excision) is the removal of the clitoris along with all or part of the labia minora. This may or may not include removing all or part of the labia majora.
Type III (Infibulation) is the act of removing the inner or outer labia and sealing the wound, leaving only a narrow opening.
Type IV refers to "all other harmful procedures to the female genitalia for non-medical purposes (piercing, scraping, cauterizing of the genital area)."
FGM often takes the form of a traditional celebration conducted by an elder or community leader. The age that women undergo the procedure varies depending on the culture, although it is most commonly performed on prepubescent girls. Certain cultures value FGM as a coming of age ritual for girls and use it to preserve a woman's virginity and faithfulness to the husband after marriage. It is also closely connected with some traditional ideals of female beauty and hygiene. FGM may or may not have religious connotations depending on the circumstances.
There are no health benefits of FGM, as it interferes with the natural functions of a woman's and girls' bodies, such as causing severe pain, shock, hemorrhage, tetanus or sepsis (bacterial infection), urine retention, open sores in the genital region and injury to nearby genital tissue, recurrent bladder and urinary tract infections, cysts, increased risk of infertility, childbirth complications and newborn deaths. Sexual problems are 1.5 more likely to occur in women who have undergone FGM, they may experience painful intercourse, have less sexual satisfaction, and be two times more likely to report a lack of sexual desire. In addition, the maternal and fetal death rate is significantly higher due to childbirth complications.
FGM can have severe negative psychological effects on women, both during and after the procedure. These can include long-term symptoms of depression, anxiety, post-traumatic stress disorder, and low self-esteem. Some women report that the procedure was carried out without their consent and knowledge, and describe feelings of fear and helplessness while it was taking place. A 2018 study found that larger quantities of the hormone cortisol were secreted in women who had undergone FGM, especially those who had experienced more severe forms of the procedure and at an early age. This marks the body's chemical response to trauma and stress and can indicate a greater risk for developing symptoms of PTSD and other trauma disorders, although there are limited studies showing a direct correlation.
The Istanbul Convention prohibits FGM (Article 38). Legislation has been introduced in certain countries to prevent FGM. A 2016 survey of 30 countries showed 24 had policies to manage and prevent FGM, although the process to provide funding, education, and resources were often inconsistent and lacking. Some countries have seen a slight decline in FGM rates, while others show little to no change.
Child and forced marriage
The practice of forcing young girls into early marriage, common in many parts of the world, is threatening their reproductive health. According to the World Health Organization:
Niger has the highest prevalence of child marriage under 18 in the world, while Bangladesh has the highest rate of marriage of girls under age 15. Practices such as bride price and dowry can contribute to child and forced marriages.
International Conference on Population and Development, 1994
The International Conference on Population and Development (ICPD) was held in Cairo, Egypt, from 5 to 13 September 1994. Delegations from 179 States took part in negotiations to finalize a Programme of Action on population and development for the next 20 years. Some 20,000 delegates from various governments, UN agencies, NGOs, and the media gathered for a discussion of a variety of population issues, including immigration, infant mortality, birth control, family planning, and the education of women.
In the ICPD Program of Action, 'reproductive health' is defined as:
This definition of the term is also echoed in the United Nations Fourth World Conference on Women, or the so-called Beijing Declaration of 1995. However, the ICPD Program of Action, even though it received the support of a large majority of UN Member States, does not enjoy the status of an international legal instrument; it is therefore not legally binding.
The Program of Action endorses a new strategy which emphasizes the numerous linkages between population and development and focuses on meeting the needs of individual women and men rather than on achieving demographic targets. The ICPD achieved consensus on four qualitative and quantitative goals for the international community, the final two of which have particular relevance for reproductive health:
Reduction of maternal mortality: A reduction of maternal mortality rates and a narrowing of disparities in maternal mortality within countries and between geographical regions, socio-economic and ethnic groups.
Access to reproductive and sexual health services including family planning: Family planning counseling, pre-natal care, safe delivery and post-natal care, prevention and appropriate treatment of infertility, prevention of abortion and the management of the consequences of abortion, treatment of reproductive tract infections, sexually transmitted infections and other reproductive health conditions; and education, counseling, as appropriate, on human sexuality, reproductive health, and responsible parenthood. Services regarding HIV/AIDS, breast cancer, infertility, delivery, hormone therapy, sex reassignment therapy, and abortion should be made available. Active discouragement of female genital mutilation (FGM).
The keys to this new approach are empowering women, providing them with more choices through expanded access to education and health services, and promoting skill development and employment. The programme advocates making family planning universally available by 2015 or sooner, as part of a broadened approach to reproductive health and rights, provides estimates of the levels of national resources and international assistance that will be required, and calls on governments to make these resources available.
Sustainable Development Goals
Half of the development goals put on by the United Nations started in 2000 to 2015 with the Millennium Development Goals (MDGs). Reproductive health was Goal 5 out of 8. To monitor the progress, the UN agreed to four indicators:
Contraceptive prevalence rates
Adolescent birth rate
Antenatal care coverage
Unmet need for family planning
Progress was slow, and according to the WHO in 2005, about 55% of women did not have sufficient antenatal care and 24% had no access to family planning services. The MDGs expired in 2015 and were replaced with a more comprehensive set of goals to cover a span of 2016–2030 with a total of 17 goals, called the Sustainable Development Goals. All 17 goals are comprehensive in nature and build off one another, but goal 3 is "To ensure healthy lives and promote wellbeing for all at all ages". Specific targets are to reduce global maternal mortality ratio to less than 70 per 100,000 live births, end preventable deaths of newborns and children, reduce the number by 50% of accidental deaths globally, strengthen the treatment and prevention programs of substance abuse and alcohol. Goal 4 emphasizes the fact that no one should be left out in providing quality education. Target 4 specifically mentions the inclusion of persons with disabilities, indigenous peoples and children in vulnerable situations. In addition, one of the targets of the Sustainable Development Goal 5 is to ensure universal access to sexual and reproductive health.
By region
North America
The CDC estimated that one in five people in the US had a sexually transmitted infection (STI) totalling near about 68 million infections in 2018. 26 million new STI in 2018. Almost half of new STI were among youth aged 15 to 24 in the US. New STIs total $16 billion in direct medical costs.
Engaging in oral sex can carry the risk of sexually transmitted infections (STIs).
AfricaHIV/AIDSHIV/AIDS in Africa is a major public health problem. The population of Sub-Saharan Africa is the worst affected region with the disease especially affecting the young female population. According to the National Library of Medicine, "Sub Saharan Africa (SSA) is occupied by 12% of the global population, but disproportionately has more than 90% of children younger than 15 years of age and 68% of adults that are living with HIV2." In Nigeria in specific, "There is early sexual maturity and considerable sexual activity between 9 and 15 years of age." HIV is also transmissible through breast milk, which proves that women infected with HIV/AIDS have to deal with more health consequences. South of the Sahara, the AIDS epidemic is the leading cause of death.
The reasons for the high spread of HIV/AIDS can be broken down into 7 main subsections: poverty, inadequate medical care, lack of prevention and education, taboo and stigma, sexual behavior, prostitution, and sexual violence against women. With a high population of individuals living in extreme poverty, condoms, HIV tests, and other forms of screening are not prioritized, leaving many individuals lacking the necessities to protect themselves from the disease. According to the International Finance Corporation, "Health care in Sub-Saharan Africa remains the worst in the world, with few countries able to spend the $34 to $40 a year per person that the World Health Organization considers the minimum for basic health care." Notably, though widespread poverty, "an astonishing 50 percent of the region's health expenditure is financed by out-of-pocket payments from individuals." This represents the lack of both affordability and accessibility surrounding the health care system in Sub-Saharan Africa. According to the United Nation, Sub-Saharan Africa struggles with the highest rate of education exclusion in the world; 60% of youth ages 15 to 17 are not in school. With this lack of education, information regarding HIV/AIDS and prevention practices are not transmitted to a number of individuals, leading to more citizens being unaware of the severity of the disease. Stigma surrounding HIV/AIDS further contributes to the high infection rate. In African villages, an individual's life is closely intertwined with their friends, families, and neighbors around them. Individuals who have HIV/AIDS are motivated to keep it a secret in fear of isolation and alienation. The extremity of this stigma is conveyed by some of the dialogue, people living with HIV are often ridiculed as "a walking corpse", referred to as "an HIV" and even called in Tanzania, "nyambizi", or submarine, which implies that an HIV-positive person is "menacing and deadly." Sexual behavior and prostitution also play a part in the increased rate of transmission of HIV/AIDS in Africa. Due to the high rates of poverty, prostitution is widespread, and sexual partners are often changing, increasing the likelihood of transmission. Africa has one of the highest rates of rape in the world, with many women getting AIDS due to raped and sexual violence by an HIV-infected offender. Similarly, gender roles within many African countries contribute to this, as "in much of sub-Saharan Africa, women are a subordinate group who are expected to become pregnant, bear children, and fulfill the sexual desires of their husbands without hesitation".Fertility rates and contraceptivesIn most African countries, the total fertility rate is very high often due to a lack of access to contraception, family planning, and practices such as forced child marriage. For instance, Niger, Angola, Mali, Burundi, Somalia and Uganda have very high fertility rates. According to the United Nations Department of Economic and Social Affairs, "Africa has the lowest rate of contraceptive use (33%) and the highest rate of unmet need for contraceptives (22%)." In Mozambique, despite efforts in improving access to modern contraceptive methods, the general fertility rate is "still high at 5.3 and the unmet need for contraceptives is also high at 26%." Among young women, the fertility rate has dramatically increased from 167 births per 1000 aged between (15–19 years) in 2011 to 194 in 2015 with a large increase in rural areas from 183 to 230. Contraceptive prevalence among (15–19 years) remains low at 14% in 2015 when compared to the national prevalence among the reproductive age group (15–49 years) at 25% in the same year.Types of contraceptivesThe copper IUD has been provided less frequently than other contraceptive methods but there have been signs of an increase in most reported provinces. The most frequently provided methods are implants and injectable progesterone, which is not as ideal as condom usage, which is still required with this method to decrease the risk of HIV. In Nigeria, specifically, people who have multiple partners are often unwilling to protect themselves with condoms. "In a study conducted in a rural community in South West Nigeria in 1993, it was found that although 94.7% of 302 candidates aged between 20 and 54 years admitted hearing about the condom, only 51.3% admitted ever using it." According to the International Family Planning Perspective, "these injectable progesterone products made up 49% of South Africa's contraceptive use and up to 90% in some provinces." Though contraceptive use is rising in African countries, discontinuation rates are also high. Weak health systems challenge Sub-Saharan African countries in expanding contraceptive outreach, promotions and service.Contraceptive accessibilityThe updated contraceptive guidelines in South Africa attempt to improve accessibility by providing special service delivery and prompting awareness for adolescents, lesbian, gay, bisexual, transgender, intersex people, disabled people, chronically ill people, women who are perimenopausal, sex workers, migrants and males. They also aim to increase access to long-acting contraceptive methods such as the copper IUD, the single rod progestogen implant combined with estrogen and progesterone injectables. Tanzanian provider perspectives also realized the biggest obstacle in maintaining healthy contraceptive care in their communities: lack of consistency. Contraceptive dispensaries found that the capability of providing service to patients was inconsistent and substandard. This resulted in unsatisfied reproductive goals, low educational attainment, miseducation about the side effects of certain contraceptives.
Accessibility has also been hindered as a result of inadequate quantities of properly trained medical personnel. According to the African Journal of Reproductive Health, "Shortage of the medical attendant...is a challenge, we are not able to attend to a big number of clients, also we do not have enough education which makes us unable to provide women with the methods they want". The majority of medical centers are staffed by people without medical training and few doctors and nurses, despite federal regulations, due to lack of resources. One center had only one person who was able to insert and remove implants, and without her, they were unable to service people who required this method of contraceptive care. Another dispensary which carried two methods of birth control shared that they sometimes run out of both materials at the same time which makes it difficult to keep up with the supply and demand chain.Social factors effect on contraceptivesUnbalanced gender dynamics, spousal dynamics, economic conditions, religious norms, cultural norms, and constraints in supply chains all contribute to contraceptive rates and usage. One instance of this is a provider who referenced harmful propaganda about the side effects of contraceptive usage. The spread of this propaganda is one of the many examples of influential people in the community, such as elders and religious leaders, discouraging proper contraceptive care/health. In some cases, influential members of the community often convince others that condoms and contraceptive pills contain microorganisms that cause cancer.
In regards to spousal and gendered dynamics, many women often have faced pressure from their spouse or family members to use avoid birth control which resulted in them using it secretly. This is also one of the many reasons women frequently preferred undetectable contraceptive methods which can lead to less effective contraceptives.Other common sexually transmitted infections in Sub-Saharan AfricaSub-Saharan Africa ranks first in STI yearly incidence compared to other world regions, reiterating the major problem that public health is in African countries. In Sub-Saharan Africa, STIs are the most common reasons that individuals seek medical care. According to the World Health Organization, every year in Africa "there are 3.5 million cases of syphilis, 15 million cases of chlamydial disease, 16 million cases of gonorrhea, and 30 million cases of trichomoniasis."Sexually transmitted infections and womenThe majority of HIV infections, risks, and other sexually transmitted infections in Sub-Saharan Africa disproportionately impact women. Women, particularly under the age of 30, account for more than half of new infections on the African continent, employing incidence rates that are often double that of their male counterparts. Not only do women contain more risk of infection, but the consequences of these diseases are often significantly worse for women, as they can affect reproductive health as well. Some consequences of bacterial STIs include "pelvic inflammatory disease, chronic pelvic pain, tubal infertility, pregnancy complications, fetal and neonatal death." HIV infection is less unbalanced in gender infections, but other STIs disproportionately affect women, "who bear 80 percent of the disability." Previously stated, women are also more susceptible to infection due to social stigma and gendered expectations. "Most women with STDs will not seek medical care at all, or will only present late for treatment, when complications have already developed, complications that have devastating physical, psychological, and social consequences, particularly for women and their children." Women of lower-income status are often the least likely subgroup to receive any sort of medical attention.More on LGBTQ+ healthIndividuals who identify as transgender often yield significantly higher rates of HIV in comparison to other subgroups. African politics and government are silent on LGBTQ+ issues in the political sphere, which translates in part to their accessibility and prioritization in healthcare. "It is possible that the invisibility of transgender people in epidemiological data from Africa is related to the criminalization of same-sex behaviour in many countries," representative of how traditional attitudes shape one's ability to participate similarly in society. Further research conducted among transgender women in South Africa shows more "health disparities and poor access to appropriate mental, sexual and reproductive health services." Still, however, there is limited data concerning transgender individuals within African countries.
Individuals identifying as part of the LGBTQ+ community, in a study conducted by BMC International Health and Human Rights, resulted all in facing some sort of discrimination by healthcare providers based on their sexual orientation and/or gender identity. Violations took four distinct forms: availability, accessibility, acceptability, and quality. Facilities in South Africa lack services for specific LGBT concerns, providers refuse to care for patients identifying within the community, and if did, articulate moral disapproval. Finally, the lack of quality and knowledge about LGBTQ+ identities and health needs contributes to disproportionate negative harms, avoiding or delaying seeking healthcare with these implications.
The workplace and reproductive health
Reproductive health can be impacted by exposures in the workplace. Both women and men who work during their reproductive years can be exposed to a variety of chemical, physical, and psychosocial hazards at work that can impact their fertility. Many women continue to work while pregnant, thus increasing the likelihood that both mother and baby could be exposed.
Routes of exposure
Harmful substances can enter a woman's body through breathing in (inhalation), contact with the skin, or swallowing (ingestion). Pregnant workers and those planning to become pregnant should be especially concerned about exposure to reproductive hazards. Some chemicals (such as alcohol) can circulate in the mother's blood, pass through the placenta, and reach the developing fetus. Other hazardous agents can affect the overall health of the woman and reduce the delivery of nutrients to the fetus. Radiation can pass directly through the mother's body to harm her eggs or the fetus. Some drugs and chemicals can also pass through a mother's body into the nursing baby through the breast milk.
Reproductive hazards do not affect every woman or every pregnancy. Whether a woman or her baby is harmed depends on how much of the hazard they are exposed to, when they are exposed, how long they are exposed, how they are exposed, and personal factors like age, stage of menstrual cycle, stage of pregnancy or when exposure occurs. For example, exposure to a hazard could block ovulation and pregnancy only at specific times of the menstrual cycle. Exposure during the first 3 months of pregnancy might cause a birth defect or a miscarriage. Exposure during the last 6 months of pregnancy could slow the baby's growth, affect its brain development, or cause premature labor.
Workplace substances that affect female workers and their pregnancies can also harm their families. Without knowing it, workers can bring home harmful substances that can affect the health of other family members—both adults and children. For example, lead brought home from the workplace on a worker's skin, hair, clothes, shoes, tool box, or car can cause lead poisoning in family members, especially young children.
Occupational reproductive hazards
A number of occupational hazards can impact reproductive health and subsequently reproductive outcomes including chemical, physical, and psychosocial hazards. Although more than 1,000 workplace chemicals have been shown to have reproductive effects on animals, most have not been studied in humans. In addition, most of the 4 million other chemical mixtures in commercial use remain untested.
Some reproductive hazards include:
Anesthetic gases
Antineoplastic (cancer treatment drugs)
Chemical disinfectants and sterilants
Certain ethylene glycol ethers such as 2-ethoxyethanol (2EE) and 2-methoxyethanol (2ME)
Carbon disulfide (CS2)
Epoxies and resins
Ethylene Oxide
Formaldehyde
Heat
Infectious agents
Lead and other heavy metals
Noise
Pesticides
Ionizing radiation
Non-ionizing radiation
Secondhand smoke
Smoke and by-products of burning
Solvents
Shift work and long working hours
Strenuous physical demands (e.g. prolonged standing, heavy lifting, bending)
Many chemicals are not evaluated for reproductive toxicity and occupational exposure limits are developed based on nonpregnant adults. Exposure levels considered safe for an adult may, or may not be safe for a fetus.
Reproductive health problems that might be caused by workplace exposures
Workplace hazards can lead to certain reproductive health problems, such as:
Reduced fertility or infertility
Erectile dysfunction
Menstrual cycle and ovulatory disorders
Women's health problems linked to sex hormone imbalance
Miscarriage
Stillbirth
Babies born too soon or too small
Birth defects
Child developmental disorders
Childhood cancers
Occupational hazards and female reproductive health
Some workplace hazards can affect reproductive health, the ability to become pregnant, and the health of unborn children. Most women can safely keep working in their job during their pregnancy. But some jobs involve exposures that are harmful to pregnant or breastfeeding women. Some female health problems that may be caused by workplace reproductive hazards include the following:
Disruption of the menstrual cycle and hormone production
High levels of physical or emotional stress or exposure to chemicals such as pesticides, polychlorinated biphenyls (PCBs), organic solvents and carbon disulfide, may disrupt the balance between the brain, pituitary gland, and ovaries. This disruption can result in an imbalance of estrogen and progesterone, and lead to changes in menstrual cycle length and regularity and ovulation. Because these sex hormones have effects throughout a woman's body, severe or long-lasting hormone imbalances may affect a woman's overall health.
Hazards that can disrupt the menstrual cycle and/or sex hormone production include:
a variety of pesticides
carbon disulfide (CS2)
polychlorinated biphenyls (PCBs)
organic solvents
jet fuel
shift work
Infertility and subfertility
About 10% to 15% of all couples are infertile or have subfertility, which means that they are unable to conceive a child after 1 year of trying to become pregnant. Many factors can affect fertility, and these factors can affect one or both partners. Damage to the woman's eggs or the man's sperm, or a change in the hormones needed to regulate the normal menstrual cycle are just a few things that can cause problems with fertility. More common causes of infertility include:
Damage to the woman's eggs
Damage to the man's sperm
Infertility can be caused by change in the hormones needed to regulate the normal menstrual cycle and uterine growth. Hazards that can reduce fertility in women include:
cancer treatment drugs, including antineoplastic drugs
lead
ionizing radiation, including x-rays and gamma rays
nitrous oxide (N2O)
Miscarriages and stillbirths
About 1 in every 6 pregnancies ends in a miscarriage—the unplanned termination of a pregnancy. Miscarriages can occur very early in pregnancy, even before the woman knows she is pregnant. Miscarriages and stillbirths occur for many reasons, such as the following:
The egg or sperm may be damaged so that the egg cannot be fertilized or cannot survive after fertilization.
A problem may exist in the hormone system needed to maintain the pregnancy.
The fetus may not have developed normally.
Physical problems may exist with the uterus or cervix.
Birth defects
A birth defect is a physical abnormality present at birth, though it may not be detected until later. About 2% to 3% of babies are born with a major birth defect. In most cases, the cause of the birth defect is unknown. The first 3 months of the pregnancy is a very sensitive time of development because the internal organs and limbs are formed during this period. Many women are not aware that they are pregnant during much of this critical period.
Low birth weight and premature birth
About 7% of babies born in the United States are born underweight or prematurely. Poor maternal nutrition, smoking, and alcohol use during pregnancy are believed to be responsible for most of these cases. Although better medical care has helped many underweight or premature babies to develop and grow normally, they are more likely than other babies to become ill or even die during their first year of life.
Developmental disorders
Sometimes the brain of the fetus does not develop normally, which leads to developmental delays or learning disabilities later in life. About 10% of children in the United States have some form of developmental disability. Such problems are often not noticeable at birth. They can be difficult to measure, may be temporary or permanent, and range from mild to severe. Developmental problems may appear as hyperactivity, short attention span, reduced learning ability, or (in severe cases) intellectual disability.
Other health problems
Even if a woman is not trying to become pregnant, her general health can be harmed by reproductive hazards that alter the production of sex hormones. Sex hormones have effects throughout a woman's body. Some workplace exposures can cause an imbalance of estrogen and progesterone levels in the blood. This disruption can increase vulnerability to:
Some cancers, such as endometrial or breast cancer
Osteoporosis
Heart disease
Tissue loss or weakening
Effects on the brain and spinal cord, including symptoms of menopause
Occupational hazards and male reproductive health
A number of workplace substances have been identified as reproductive hazards for men such as:
Lead
Dibromochloropropane
Carbaryl (sevin)
Toluenediamine and dinitrotoluene
Ethylene dibromide
Plastic production (styrene and acetone)
Ethylene glycol monoethyl ether
Welding
Perchloroethylene
Mercury vapor
Heat
Military radar
High levels of kepone
High levels of bromine vapor
High levels of radiation
Carbon disulfide
2,4-dichlorophenoxy acetic acid (2,4-D)
Exposure to occupational hazards can impact:
Number of sperm. Some reproductive hazards can stop or slow the actual production of sperm. This means that there will be fewer sperm present to fertilize an egg; if no sperm are produced, the man is sterile. If the hazard prevents sperm from being made, sterility is permanent.
Sperm shape. Reproductive hazards may cause the shape of sperm cells to be different. These sperm often have trouble swimming or lack the ability to fertilize the egg.
Sperm transfer. Hazardous chemicals may collect in the epididymis, seminal vesicles, or prostate. These chemicals may kill the sperm, change the way in which they swim, or attach to the sperm and be carried to the egg or the unborn child.
Sexual performance. Changes in amounts of hormones can affect sexual performance. Some chemicals, like alcohol, may also affect the ability to achieve erections, whereas others may affect the sex drive. Several drugs (both legal and illegal) have effects on sexual performance, but little is known about the effects of workplace hazards.
Sperm chromosomes. Reproductive hazards can affect the chromosomes found in sperm. The sperm and egg each contribute 23 chromosomes at fertilization. The DNA stored in these chromosomes determines what someone will look like and their our bodies will function. Radiation or chemicals may cause changes or breaks in the DNA. If the sperm's DNA is damaged, it may not be able to fertilize an egg; or if it does fertilize an egg, it may affect the development of the fetus. Some cancer treatment drugs are known to cause such damage. However, little is known about the effects of workplace hazards on sperm chromosomes.
Pregnancy'''. If a damaged sperm does fertilize an egg, the egg might not develop properly, causing a miscarriage or a possible health problem in the baby. If a reproductive hazard is carried in the semen, the fetus might be exposed within the uterus, possibly leading to problems with the pregnancy or with the health of the baby after it is born.
| Biology and health sciences | Fields of medicine | null |
11737869 | https://en.wikipedia.org/wiki/Snow%20pea | Snow pea | The snow pea is an edible-pod pea with flat pods and thin pod walls. It is eaten whole, with both the seeds and the pod, while still unripened.
Names
The common name snow pea seems to be a misnomer as the planting season of this pea is no earlier than that of other peas. Another common name, Chinese pea, is probably related to its prominence in Chinese dishes served in the West. It is called mangetout in the United Kingdom and Ireland (from the French for "eat-all" and pronounced monge-too; /mɒnʒtuː/).
Snow peas and snap peas both belong to Macrocarpon Group, a cultivar group based on the variety Pisum sativum var. macrocarpum Ser. named in 1825. It was described as having very compressed non-leathery edible pods in the original publication.
The scientific name Pisum sativum var. saccharatum Ser. is often misused for snow peas. The variety under this name was described as having sub-leathery and compressed-testes pods and the French name petit pois. The description is inconsistent with the appearance of snow peas, and therefore botanists have replaced this name with Pisum sativum var. macrocarpum. Austrian scientist and monk Gregor Mendel used peas which he called Pisum saccharatum in his famous experiments demonstrating the heritable nature of specific traits, and this Latin name might not refer to the same varieties identified with modern snow peas.
Composition
Nutrition
Uses
Culinary
Snow peas, along with sugar snap peas and unlike field and garden peas, are notable for having edible pods that lack inedible fiber (in the form of "parchment", a fibrous layer found in the inner pod rich in lignin) in the pod walls. Snow peas have the thinner walls of the two edible pod variants. Two recessive genes known as p and v are responsible for this trait. p is responsible for reducing the sclerenchymatous membrane on the inner pod wall, while v reduces pod wall thickness (n is a gene that thickens pod walls in snap peas).
Pea shoots () are the stems and leaves of the immature plant, used as a vegetable in Chinese cooking. They are commonly stir-fried with garlic and sometimes combined with crab or other shellfish.
Nitrogen fixers
As with most legumes, snow peas host beneficial bacteria, rhizobia, in their root nodules, which fix nitrogen in the soil—this is called a mutualistic relationship—and are therefore a useful companion plant, especially useful to grow intercropped with green, leafy vegetables that benefit from high nitrogen content in their soil.
Cultivation
Snow peas can be grown in open fields during cool seasons and can thus be cultivated during winter and spring seasons.
Storage
Storage of the pea with films of polymethylpentene at a temperature of and controlled atmosphere with a concentration of oxygen and carbon dioxide of 5 kPa augments the shelf life, internal and external characteristics of the plant.
Gallery
| Biology and health sciences | Pulses | Plants |
11738757 | https://en.wikipedia.org/wiki/Peak%20demand | Peak demand | Peak demand on an electrical grid is the highest electrical power demand that has occurred over a specified time period (Gönen 2008). Peak demand is typically characterized as annual, daily or seasonal and has the unit of power.
Peak demand, peak load or on-peak are terms used in energy demand management describing a period in which electrical power is expected to be provided for a sustained period at a significantly higher than average supply level. Peak demand fluctuations may occur on daily, monthly, seasonal and yearly cycles. For an electric utility company, the actual point of peak demand is a single half-hour or hourly period which represents the highest point of customer consumption of electricity. At this time there is a combination of office, domestic demand and at some times of the year, the fall of darkness.
Some utilities will charge customers based on their individual peak demand. The highest demand during each month or even a single 15 to 30 minute period of highest use in the previous year may be used to calculate charges. The renewable energy transition will include considerations for peak demand.
Economic growth of the state is inversely associated with peak load.
Demand Tariff
Electricity network is built to deal with the highest possible peak demand otherwise blackout may happen. In Australia, demand tariff has three components: peak demand charge, energy charge and daily connection charge. For example, for large customers (commercial, industrial or mixed of commercial/residential), the peak demand charge is based on the highest 30 minutes electricity consumption in a month; the energy charge is based on a month electricity consumption. This type of demand tariff is gradually introduced to residential households and will be rolled out by 2020 in Queensland Australia. How to manage electricity bills under demand tariff can be challenging. The key solutions involve improving building efficiency and managing the operational settings of large power appliances.
Time of Peak Demand
Peak Demand depends on the demography, the economy, the weather, the climate, the season, the day of the week and other factors. In industrialised regions of China or Germany, the peak demands mostly occur in day time. However, in more service based economy such as Australia, the daily peak demands often occur in the late afternoon to early evening time (e.g. 4pm to 8pm). Residential and commercial electricity demand contributes a lot to this type of network peak demand.
Off-peak
Peak demand is considered to be the opposite to off-peak hours when power demand is usually low. There are off-peak time-of-use rates. Sometimes, there are 3 time-of-use zones: peak, shoulder and offpeak. Shoulder is often the time between peak and offpeak in weekdays. Weekends are often just peak and offpeak in terms of managing electricity loads for the network.
Response
Peak demand may exceed the maximum supply levels that the electrical power industry can generate, resulting in power outages and load shedding. This often occurs during heat waves when use of air conditioners and powered fans raises the rate of energy consumption significantly. During a shortage authorities may request the public to curtail their energy use and shift it to a non-peak period.
Power stations
Power stations specifically constructed for providing power to electrical grids for peak demand are called peaking power plants or 'peakers'. In general, Natural gas fueled power stations can be fired up rapidly and are therefore often utilized at peak demand times. Combined cycle power plants can frequently provide power for peak demand, as well as run efficiently for baseload power.
Hydroelectric power and pumped storage type dams such as Carters Dam in the U.S. state of Georgia help to meet peak demand as well.
The chances that a wind farm will be unable to meet peak demand are greater than for a fossil-fueled power station, due to the ability to store liquid fuels for use during peak demand.
Solar power's peak output often naturally coincides with daytime peaks of usage due to air conditioning.
| Technology | Concepts | null |
13256745 | https://en.wikipedia.org/wiki/Planetary%20nebula%20luminosity%20function | Planetary nebula luminosity function | Planetary nebula luminosity function (PNLF) is a secondary distance indicator used in astronomy. It makes use of the [O III] λ5007 forbidden line found in all planetary nebula (PNe) which are members of the old stellar populations (Population II).
It can be used to determine distances to both spiral and elliptical galaxies despite their completely different stellar populations and is part of the Extragalactic Distance Scale.
Procedure
The distance estimate to a galaxy using the PNLF requires discovery of such an object in the target galaxy that is visible at λ5007 but not when the entire spectrum is considered. These points are candidate PNe, however, there are three other types of objects that would also exhibit such an emission line that must be filtered out: HII regions, supernova remnants, and Lyα galaxies. After the PNe are determined, to estimate a distance one must measure their monochromatic [O III] λ5007 luminosity. What remains is a statistical sample of
PNe. The observed luminosity function is then fitted to some standard law.
Finally, one must estimate the foreground interstellar extinction. The two sources of extinction, are from within the Milky Way and the internal extinction of the target galaxy. The first is well known and can be taken from sources such as reddening maps computed from H I measurements and galaxy counts or from IRAS and DIRBE satellite experiments. The later type of extinction, occurs only in target galaxies which are either late type spiral or irregular. However, this extinction is difficult to measure. In the Milky Way, the scale height of PNe is much bigger than that of the dust. Observational data and models support that this holds true for other galaxies, that the bright edge of the PNLF is primarily due to PNe in front of the dust layer. The data and models support a less than 0.05 apparent magnitude internal extinction of a galaxy's PNe.
Physics behind process
The PNLF method is unbiased by metallicity. This is because oxygen is a primary nebular coolant; any drop in its concentration raises the plasma's electron temperature and raises the amount of collisional excitations per ion. This compensates for having a smaller number of emitting ions in the PNe resulting in little change in the λ5007 emissions . Consequently, a reduction in oxygen density only lowers the emergent [O III] λ5007 emission line intensity by approximately the square root of the difference in abundance. At the same time, the PNe's core responds to metallicity the opposite way. In the case where the metallicity of the progenitor star is smaller, the PNe's central star will be a bit more massive and its illuminating ultraviolet flux will be a bit greater. This added energy almost precisely accounts for the decreased emissions of the PNe. Consequently, the total [O III] λ5007 luminosity that is produced by a PNe is practically
uncorrelated to metallicity. This beneficial negation is in agreement with more precise models of PNe evolution. Only in extremely metal-poor PNe does the brightness of the PNLF cutoff dim by more than a small percentage.
The relative independence of the PNLF cutoff with respect to population age is harder to understand. The [O III] λ5007 flux of a PNe directly correlates to the brightness of its central star. Further, the brightness of its central star directly correlates to its mass and the central star's mass directly varies in relation to its progenitor's mass. However, by observation, it is demonstrated that reduced brightness does not happen.
| Physical sciences | Basics | Astronomy |
569684 | https://en.wikipedia.org/wiki/Utahraptor | Utahraptor | Utahraptor (meaning "Utah's predator") is a genus of large dromaeosaurid (a group of feathered carnivorous theropods) dinosaur that lived during the Early Cretaceous period from around 135 to 130 million years ago in what is now the United States. The genus was described in 1993 by American paleontologist James Kirkland and colleagues with the type species Utahraptor ostrommaysi, based on fossils that had been unearthed earlier from the Cedar Mountain Formation of Utah. Later, many additional specimens were described including those from the skull and postcranium in addition to those of younger individuals.
The genus contains a single species, Utahraptor ostrommaysi. It is the largest-known member of the family Dromaeosauridae, measuring about long and typically weighing up to . As a heavily built, ground-dwelling, bipedal carnivore, its large size and variety of unique features have earned it attention in both pop culture and the scientific community. The jaws of Utahraptor were lined with small, serrated teeth that were used in conjunction with a large "killing claw" on its second toe to dispatch its prey. Its skull was boxy and elongated, akin to other dromaeosaurids like Dromaeosaurus and Velociraptor.
Being a carnivore, Utahraptor was adapted to hunt the other animals of the Cedar Mountain Formation ecosystem such as ankylosaurs and iguanodonts. Evidence from the leg physiology supports the idea of Utahraptor being an ambush predator, in contrast to other dromaeosaurs that were pursuit predators. Fossil remains of several individuals of various ages have been found together, suggesting that Utahraptor was gregarious (social) and practiced degrees of post nestling care.
Discovery and naming
The first specimens of Utahraptor were found in 1975 by Jim Jensen in the Dalton Wells Quarry of Utah, near the town of Moab, but did not receive much attention. After a find of a large claw by Carl Limone in October 1991, James Kirkland, Robert Gaston and Donald Burge uncovered further remains of Utahraptor in 1991 in the Gaston Quarry in Grand County, Utah, within the Yellow Cat and Poison Strip members of the Cedar Mountain Formation. The holotype of Utahraptor, CEUM 184v.86, consists of a second pedal ungual, with potentially assigned elements from other specimens: pedal ungual CEUM 184v.294, tibia CEUM 184v.260 and premaxilla CEUM 184v.400. The holotype is housed in the paleontology collections of the Prehistoric Museum at Utah State University Eastern. Brigham Young University, the depository of Jensen's finds, currently houses the largest collection of Utahraptor fossils.
The type species, Utahraptor ostrommaysi, was named by Kirkland, Gaston and Burge in June 1993. The genus name Utahraptor is in reference to Utah, where the remains were found. The specific name, ostrommaysi, is in honor to John Ostrom for his investigations on Deinonychus and its relationships to birds, as well as Chris Mays, who helped in the research of Utahraptor by founding Dinamation. From his description, Kirkland stated the meaning of genus name to be "Utah's predator," but the Latin word raptor translates to 'robber' or 'plunderer', not 'predator'. Earlier, it had been intended to name the species "U. spielbergi" after film director Steven Spielberg, in exchange for him funding paleontological research, but no agreement could be reached on the amount of financial assistance.
In 2000, the specific name was emended by George Olshevsky to the plural genitive ostrommaysorum. However, Thiago Vernaschi V. Costa and Normand David in 2019 criticized the use of the species name U. ostrommaysorum, since it has no clear justification or explanation. Although this spelling has been largely used by other authors, the genus Utahraptor was originally coined with the type species U. ostrommaysi and, given that the International Code of Zoological Nomenclature offers no provision for forming a genitive form from two persons with different names, Costa and David conclude that the original spelling ostrommaysi has to be regarded as an arbitrary combination of letters and not a correctly formed genitive form. Under this reasoning, ostrommaysorum has no valid use and the original spelling ostrommaysi does not need to be emended. Other alternative and also invalid spellings were used in scientific literature, such as ostromaysi, ostromaysorum, ostromayssorum, ostromayorum and ostrommaysori.
Some elements were wrongly referred to the genus. The lacrimal bone of the specimen CEUM 184v.83 turned out to be a postorbital from the ankylosaur Gastonia. Britt et al. also suggested that the previously identified manual unguals of the specimens M184v.294, BYU 9438 and BYU 13068 are indeed pedal unguals. This suggestion was confirmed by Senter in 2007.
Description
Utahraptor was one of, if not the largest and heaviest of all dromaeosaurids, with the largest assigned specimen BYUVP 15465 having a femoral length of .
Utahraptor is estimated to have reached in length and somewhat less than , comparable in weight to a polar bear. Some authors estimated that it weighed up to . In 2024, the body mass of BYUVP 2536 and BYUVP 1833 were estimated around respectively, though BYUVP 7510-18078 was estimated to have weighed .
Although feathers have never been found in association with Utahraptor specimens, there is strong phylogenetic evidence suggesting that all dromaeosaurids had them. The feathered genus Microraptor is one of the oldest-known dromaeosaurids and is phylogenetically more primitive than Utahraptor. Since Microraptor and other dromaeosaurids possessed feathers, it is reasonable to assume that this trait was present in all of Dromaeosauridae. Feathers were very unlikely to have evolved more than once, so assuming that any given dromaeosaurid, such as Utahraptor, lacked feathers would require positive evidence that they did not have them. So far, there is nothing to suggest that feathers were lost in larger, more derived species of dromaeosaurs. The presence of quill knobs in Dakotaraptor evidenced that even larger dromaeosaurids had feathers.
According to Kirkland et al. in 1993, Utahraptor can be recognized by a few special autapomorphies. The claws on its hand are more specialized as cutting blades than in other dromaeosaurids. It has a lacrimal bone with distinctly parallel mesial and outer sides that gives it an elongate subrectangular appearance in top view and it has a base of the nasal opening on the premaxilla parallel to the premaxillary tooth row. In the revised diagnosis conducted by Turner et al. in 2012, Utahraptor differs from other dromaeosaurids in having an elongate nasal process of the premaxilla, a distal end of metatarsal III that is smooth, not ginglymoid, an L-shaped quadratojugal without a posterior process, the presence of a well-developed notch between the lesser trochanter and greater trochanter, and dorsal vertebrae that lack pleurocoels. Like other dromaeosaurids, Utahraptor had a large curved claw on the second toe of each foot. The second pedal ungual is preserved with a outside curve length and is estimated to reach in restoration.
Classification
Utahraptor is a member of the family Dromaeosauridae, a clade of theropod dinosaurs commonly known as "raptors". Utahraptor is the largest known genus in the family and belongs to the same clade of other notable dinosaurs such as Velociraptor, Deinonychus, or Dromaeosaurus. It is classified in the subfamily Dromaeosaurinae, which is found in the clade Eudromaeosauria.
In 2015, Utahraptor was found to be closely related to the smaller Dromaeosaurus and the giant Mongolian and North American dromaeosaurid genera Achillobator and Dakotaraptor:
The cladogram below is the result of a cladistic analysis conducted by Cau et al. in 2017.
Paleobiology
Predatory behavior
Kirkland et al. noted that given the huge size of Utahraptor, it was not as fast as Deinonychus and Velociraptor; instead, it would have had a similar speed to the contemporary iguanodonts, and was faster than sauropods. Additionally, the thickness of the tibia indicates that the animal possessed a significant leg force in order to kill prey. It was also suggested that lighter dromaeosaurids such as Velociraptor and Deinonychus relied on their hand claws to handle prey and retain balance while kicking it; in contrast to this, the heavily built Utahraptor may have been able to deliver kicks without the risk of losing balance, freeing the hands and using them to dispatch prey.
According to Gregory S. Paul, Utahraptor was not particularly fast and would have been an ambush hunter that preyed on large dinosaurs such as the contemporary iguanodonts and therizinosaurs. Its robust build and large sickle claw indicate it was well suited to hunting such prey. Like other dromaeosaurine dromaeosaurids, it may have also relied heavily on its jaws to dispatch prey—more so than other types of dromaeosaurids, such as velociraptorines.
Social behavior
In 2001, Kirkland et al. pursued a graduate student's discovery of a bone protruding from a 9-ton fossil block of sandstone in eastern Utah. It was determined to contain the bones of at least seven individuals, including an adult measuring about , four juveniles, and a hatchling about long. Also fossilized with the Utahraptor pack are the remains of at least one possible iguanodont. Kirkland speculated that the Utahraptor pack attempted to scavenge carrion or attack helpless prey mired in quicksand, and were themselves mired in the attempt to feed on the herbivore. Similar sites such as the Cleveland-Lloyd Quarry and California's La Brea Tar Pits house such predator traps. Examination of the fossils are ongoing after a decade of excavation, but if Kirkland is correct, it may be one of the best-preserved predator traps ever discovered. The fossils may further reveal aspects into the behavior of Utahraptor, such as whether it might have hunted in groups like Deinonychus was believed to have done. Whether all the Utahraptor individuals were mired simultaneously or were drawn in, one-by-one is unclear. Further examination of the block suggests that the number of Utahraptor remains may be double the amount previously assumed.
While dinosaur behavior can only be theorized, it was later discovered in 2020 that Deinonychus may not have practiced mammal-like pack hunting, based on differing dietary preferences in adults and juveniles. Despite this, the authors stated that gregariousness was still possible for Deinonychus and the discovery of Utahraptor in the mud-trap implies it exhibited a degree of post nestling care and gregariousness.
Paleoenvironment
Utahraptor lived in the lower part of the Cedar Mountain Formation, a bed known as the Yellow Cat Member. According to the authors of its description, Utahraptor had an important ecological role as a major carnivore of the paleofauna of the present-day Arches region during the Early Cretaceous, and could probably attack prey larger than itself. Group hunting of individuals of at least and , if proven, could have killed prey of a weight of . Additionally, sauropods ranging around may have been an important part of its diet. The paleontologist Thomas R. Holtz estimated that Utahraptor existed between 130 million and 125 million years ago. In multiple occasions, the Yellow Cat Member has been dated to Barremian-Aptian ages. Sames and Schudack (2010) proposed a reassignment of the estimated age, compromising Berriasian to Valanginian stages; however, this interpretation was not followed by most authors. Using advanced methods of radiometric and palynological dating, Joeckel et al. (2019) concluded that the Yellow Cat Member is indeed older than previous estimations. The deposition occurred between 139 ± 1.3 million to 134.6 ± 1.7 million years ago, or, Berriasian to Late Valanginian stages. Based on the presence of new palynoflora, Middle Berriasian–Early Hauterivian ages were provisionally assigned. However, the Yellow Cat Member is divided into distinct "lower" and "upper" layers, and Utahraptor fossils are only currently known within the upper Yellow Cat Member.
Utahraptor was unearthed from the Yellow Cat Member, which during the Berriasian to Late Valanginian was a semiarid area with floodplain prairies, riverine forests, and open woodlands predominated by conifers (Pinophyta), ferns (Polypodiopsida), hornworts (Anthocerotophyta) and other vascular plants. During the description of Mierasaurus, it was interpreted that there was also a waterlogged bog-like environment. There is believed to have been a short wet season. This is supported by the presence of charred spores and other carbonized plant debris in the pollen maceral that indicate the occurrence of ancient wildfires ignited during periods of low precipitation.
Paleofauna that were contemporaneous with the dromaeosaurid in the upper Yellow Cat Member included numerous dinosaurs, such as the medium-sized iguanodonts Hippodraco and Cedrorestes, the smaller theropods Martharaptor and Nedcolbertia, the nodosaurid Gastonia, and the sauropods Cedarosaurus and Moabosaurus. The only known mammal from the Upper Yellow Cat Member is Cifelliodon.
Other non-dinosaur or avian taxa known from the Member include the fish Ceratodus and Semionotus, the turtles Glyptops and Trinitichelys, Aquatilavipes (fossilized bird tracks), the rhynchocephalian Toxolophosaurus, and the indeterminate remains of hybodontid and polyacrodontid sharks.
Additional paleofauna was recovered, most of it being unnamed and/or indeterminate, including an isolated mesoeucrocodylian skull that measures in length. A neochoristodere unearthed from the Upper Yellow Cat Member, represented by a partial left femur, shows that aquatic paleofauna was present and diverse during the Early Cretaceous of the Cedar Mountain Formation. A large sail-backed iguanodont represented by large vertebrae and fragmentary remains, and an indeterminate eudromaeosaur known from a caudal vertebra and fragmented tail (UMNH VP 20209) were also present.
Cultural significance
Raptor Red was published in 1995, and features the fictionalized story of a female Utahraptor. Written by paleontologist Robert T. Bakker, it was positively regarded by mainstream reviewers, though updates to the science have rendered some of the story line facts presented untrue and the paleontology community was critical of fossil record inaccuracies. Bakker's anthropomorphosis of the titular Red was particularly praised.
In 2018, it was proposed by a 10-year-old elementary school student, Kenyon Roberts, that Utahraptor be the Utah state dinosaur, an act that was approved by the Senate. Initially Utahraptor would have replaced another dinosaur, Allosaurus, as the state's official fossil, but it was decided that Utahraptor would be another symbol of the state. In 2021, Steve Eliason successfully created a proposal for Utahraptor State Park where the block was discovered, proposed by the same Utah student, Kenyon Roberts. It was approved by the state House.
| Biology and health sciences | Theropods | Animals |
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