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24603	https://en.wikipedia.org/wiki/Proteasome	Proteasome	Proteasomes are protein complexes which degrade ubiquitin-tagged proteins by proteolysis, a chemical reaction that breaks peptide bonds. Enzymes that help such reactions are called proteases. Proteasomes are part of a major mechanism by which cells regulate the concentration of particular proteins and degrade misfolded proteins. Proteins are tagged for degradation with a small protein called ubiquitin. The tagging reaction is catalyzed by enzymes called ubiquitin ligases. Once a protein is tagged with a single ubiquitin molecule, this is a signal to other ligases to attach additional ubiquitin molecules. The result is a polyubiquitin chain that is bound by the proteasome, allowing it to degrade the tagged protein. The degradation process yields peptides of about seven to eight amino acids long, which can then be further degraded into shorter amino acid sequences and used in synthesizing new proteins. Proteasomes are found inside all eukaryotes and archaea, and in some bacteria. In eukaryotes, proteasomes are located both in the nucleus and in the cytoplasm. In structure, the proteasome is a cylindrical complex containing a "core" of four stacked rings forming a central pore. Each ring is composed of seven individual proteins. The inner two rings are made of seven β subunits that contain three to seven protease active sites. These sites are located on the interior surface of the rings, so that the target protein must enter the central pore before it is degraded. The outer two rings each contain seven α subunits whose function is to maintain a "gate" through which proteins enter the barrel. These α subunits are controlled by binding to "cap" structures or regulatory particles that recognize polyubiquitin tags attached to protein substrates and initiate the degradation process. The overall system of ubiquitination and proteasomal degradation is known as the ubiquitin–proteasome system. The proteasomal degradation pathway is essential for many cellular processes, including the cell cycle, the regulation of gene expression, and responses to oxidative stress. The importance of proteolytic degradation inside cells and the role of ubiquitin in proteolytic pathways was acknowledged in the award of the 2004 Nobel Prize in Chemistry to Aaron Ciechanover, Avram Hershko and Irwin Rose. Discovery Before the discovery of the ubiquitin–proteasome system, protein degradation in cells was thought to rely mainly on lysosomes, membrane-bound organelles with acidic and protease-filled interiors that can degrade and then recycle exogenous proteins and aged or damaged organelles. However, work by Joseph Etlinger and Alfred L. Goldberg in 1977 on ATP-dependent protein degradation in reticulocytes, which lack lysosomes, suggested the presence of a second intracellular degradation mechanism. This was shown in 1978 to be composed of several distinct protein chains, a novelty among proteases at the time. Later work on modification of histones led to the identification of an unexpected covalent modification of the histone protein by a bond between a lysine side chain of the histone and the C-terminal glycine residue of ubiquitin, a protein that had no known function. It was then discovered that a previously identified protein associated with proteolytic degradation, known as ATP-dependent proteolysis factor 1 (APF-1), was the same protein as ubiquitin. The proteolytic activities of this system were isolated as a multi-protein complex originally called the multi-catalytic proteinase complex by Sherwin Wilk and Marion Orlowski. Later, the ATP-dependent proteolytic complex that was responsible for ubiquitin-dependent protein degradation was discovered and was called the 26S proteasome. Much of the early work leading up to the discovery of the ubiquitin proteasome system occurred in the late 1970s and early 1980s at the Technion in the laboratory of Avram Hershko, where Aaron Ciechanover worked as a graduate student. Hershko's year-long sabbatical in the laboratory of Irwin Rose at the Fox Chase Cancer Center provided key conceptual insights, though Rose later downplayed his role in the discovery. The three shared the 2004 Nobel Prize in Chemistry for their work in discovering this system. Although electron microscopy data revealing the stacked-ring structure of the proteasome became available in the mid-1980s, the first structure of the proteasome core particle was not solved by X-ray crystallography until 1994. In 2018, the first atomic structures of the human 26S proteasome holoenzyme in complex with a polyubiquitylated protein substrate were solved by cryogenic electron microscopy, revealing mechanisms by which the substrate is recognized, deubiquitylated, unfolded and degraded by the human 26S proteasome. Structure and organization The proteasome subcomponents are often referred to by their Svedberg sedimentation coefficient (denoted S). The proteasome most exclusively used in mammals is the cytosolic 26S proteasome, which is about 2000 kilodaltons (kDa) in molecular mass containing one 20S protein subunit and two 19S regulatory cap subunits. The core is hollow and provides an enclosed cavity in which proteins are degraded; openings at the two ends of the core allow the target protein to enter. Each end of the core particle associates with a 19S regulatory subunit that contains multiple ATPase active sites and ubiquitin binding sites; it is this structure that recognizes polyubiquitinated proteins and transfers them to the catalytic core. An alternative form of regulatory subunit called the 11S particle can associate with the core in essentially the same manner as the 19S particle; the 11S may play a role in degradation of foreign peptides such as those produced after infection by a virus. 20S core particle The number and diversity of subunits contained in the 20S core particle depends on the organism; the number of distinct and specialized subunits is larger in multicellular than unicellular organisms and larger in eukaryotes than in prokaryotes. All 20S particles consist of four stacked heptameric ring structures that are themselves composed of two different types of subunits; α subunits are structural in nature, whereas β subunits are predominantly catalytic. The α subunits are pseudoenzymes homologous to β subunits. They are assembled with their N-termini adjacent to that of the β subunits. The outer two rings in the stack consist of seven α subunits each, which serve as docking domains for the regulatory particles and the alpha subunits N-termini () form a gate that blocks unregulated access of substrates to the interior cavity. The inner two rings each consist of seven β subunits and in their N-termini contain the protease active sites that perform the proteolysis reactions. Three distinct catalytic activities were identified in the purified complex: chymotrypsin-like, trypsin-like and peptidylglutamyl-peptide hydrolyzing. The size of the proteasome is relatively conserved and is about 150 angstroms (Å) by 115 Å. The interior chamber is at most 53 Å wide, though the entrance can be as narrow as 13 Å, suggesting that substrate proteins must be at least partially unfolded to enter. In archaea such as Thermoplasma acidophilum, all the α and all the β subunits are identical, whereas eukaryotic proteasomes such as those in yeast contain seven distinct types of each subunit. In mammals, the β1, β2, and β5 subunits are catalytic; although they share a common mechanism, they have three distinct substrate specificities considered chymotrypsin-like, trypsin-like, and peptidyl-glutamyl peptide-hydrolyzing (PHGH). Alternative β forms denoted β1i, β2i, and β5i can be expressed in hematopoietic cells in response to exposure to pro-inflammatory signals such as cytokines, in particular, interferon gamma. The proteasome assembled with these alternative subunits is known as the immunoproteasome, whose substrate specificity is altered relative to the normal proteasome. Recently an alternative proteasome was identified in human cells that lack the α3 core subunit. These proteasomes (known as the α4-α4 proteasomes) instead form 20S core particles containing an additional α4 subunit in place of the missing α3 subunit. These alternative 'α4-α4' proteasomes have been known previously to exist in yeast. Although the precise function of these proteasome isoforms is still largely unknown, cells expressing these proteasomes show enhanced resistance to toxicity induced by metallic ions such as cadmium. 19S regulatory particle The 19S particle in eukaryotes consists of 19 individual proteins and is divisible into two subassemblies, a 9-subunit base that binds directly to the α ring of the 20S core particle, and a 10-subunit lid. Six of the nine base proteins are ATPase subunits from the AAA Family, and an evolutionary homolog of these ATPases exists in archaea, called PAN (proteasome-activating nucleotidase). The association of the 19S and 20S particles requires the binding of ATP to the 19S ATPase subunits, and ATP hydrolysis is required for the assembled complex to degrade folded and ubiquitinated proteins. Note that only the step of substrate unfolding requires energy from ATP hydrolysis, while ATP-binding alone can support all the other steps required for protein degradation (e.g., complex assembly, gate opening, translocation, and proteolysis). In fact, ATP binding to the ATPases by itself supports the rapid degradation of unfolded proteins. However, while ATP hydrolysis is required for unfolding only, it is not yet clear whether this energy may be used in the coupling of some of these steps. In 2012, two independent efforts have elucidated the molecular architecture of the 26S proteasome by single particle electron microscopy. In 2016, three independent efforts have determined the first near-atomic resolution structure of the human 26S proteasome in the absence of substrates by cryo-EM. In 2018, a major effort has elucidated the detailed mechanisms of deubiquitylation, initiation of translocation and processive unfolding of substrates by determining seven atomic structures of substrate-engaged 26S proteasome simultaneously. In the heart of the 19S, directly adjacent to the 20S, are the AAA-ATPases (AAA proteins) that assemble to a heterohexameric ring of the order Rpt1/Rpt2/Rpt6/Rpt3/Rpt4/Rpt5. This ring is a trimer of dimers: Rpt1/Rpt2, Rpt6/Rpt3, and Rpt4/Rpt5 dimerize via their N-terminal coiled-coils. These coiled-coils protrude from the hexameric ring. The largest regulatory particle non-ATPases Rpn1 and Rpn2 bind to the tips of Rpt1/2 and Rpt6/3, respectively. The ubiquitin receptor Rpn13 binds to Rpn2 and completes the base sub-complex. The lid covers one half of the AAA-ATPase hexamer (Rpt6/Rpt3/Rpt4) and, unexpectedly, directly contacts the 20S via Rpn6 and to lesser extent Rpn5. The subunits Rpn9, Rpn5, Rpn6, Rpn7, Rpn3, and Rpn12, which are structurally related among themselves and to subunits of the COP9 complex and eIF3 (hence called PCI subunits) assemble to a horseshoe-like structure enclosing the Rpn8/Rpn11 heterodimer. Rpn11, the deubiquitinating enzyme, is placed at the mouth of the AAA-ATPase hexamer, ideally positioned to remove ubiquitin moieties immediately before translocation of substrates into the 20S. The second ubiquitin receptor identified to date, Rpn10, is positioned at the periphery of the lid, near subunits Rpn8 and Rpn9. Conformational changes of 19S The 19S regulatory particle within the 26S proteasome holoenzyme has been observed in six strongly differing conformational states in the absence of substrates to date. A hallmark of the AAA-ATPase configuration in this predominant low-energy state is a staircase- or lockwasher-like arrangement of the AAA-domains. In the presence of ATP but absence of substrate three alternative, less abundant conformations of the 19S are adopted primarily differing in the positioning of the lid with respect to the AAA-ATPase module. In the presence of ATP-γS or a substrate, considerably more conformations have been observed displaying dramatic structural changes of the AAA-ATPase module. Some of the substrate-bound conformations bear high similarity to the substrate-free ones, but they are not entirely identical, particularly in the AAA-ATPase module. Prior to the 26S assembly, the 19S regulatory particle in a free form has also been observed in seven conformational states. Notably, all these conformers are somewhat different and present distinct features. Thus, the 19S regulatory particle can sample at least 20 conformational states under different physiological conditions. Regulation of the 20S by the 19S The 19S regulatory particle is responsible for stimulating the 20S to degrade proteins. A primary function of the 19S regulatory ATPases is to open the gate in the 20S that blocks the entry of substrates into the degradation chamber. The mechanism by which the proteasomal ATPase open this gate has been recently elucidated. 20S gate opening, and thus substrate degradation, requires the C-termini of the proteasomal ATPases, which contains a specific motif (i.e., HbYX motif). The ATPases C-termini bind into pockets in the top of the 20S, and tether the ATPase complex to the 20S proteolytic complex, thus joining the substrate unfolding equipment with the 20S degradation machinery. Binding of these C-termini into these 20S pockets by themselves stimulates opening of the gate in the 20S in much the same way that a "key-in-a-lock" opens a door. The precise mechanism by which this "key-in-a-lock" mechanism functions has been structurally elucidated in the context of human 26S proteasome at near-atomic resolution, suggesting that the insertion of five C-termini of ATPase subunits Rpt1/2/3/5/6 into the 20S surface pockets are required to fully open the 20S gate. Other regulatory particles 20S proteasomes can also associate with a second type of regulatory particle, the 11S regulatory particle, a heptameric structure that does not contain any ATPases and can promote the degradation of short peptides but not of complete proteins. It is presumed that this is because the complex cannot unfold larger substrates. This structure is also known as PA28, REG, or PA26. The mechanisms by which it binds to the core particle through the C-terminal tails of its subunits and induces α-ring conformational changes to open the 20S gate suggest a similar mechanism for the 19S particle. The expression of the 11S particle is induced by interferon gamma and is responsible, in conjunction with the immunoproteasome β subunits, for the generation of peptides that bind to the major histocompatibility complex. Yet another type of non-ATPase regulatory particle is the Blm10 (yeast) or PA200/PSME4 (human). It opens only one α subunit in the 20S gate and itself folds into a dome with a very small pore over it. Assembly The assembly of the proteasome is a complex process due to the number of subunits that must associate to form an active complex. The β subunits are synthesized with N-terminal "propeptides" that are post-translationally modified during the assembly of the 20S particle to expose the proteolytic active site. The 20S particle is assembled from two half-proteasomes, each of which consists of a seven-membered pro-β ring attached to a seven-membered α ring. The association of the β rings of the two half-proteasomes triggers threonine-dependent autolysis of the propeptides to expose the active site. These β interactions are mediated mainly by salt bridges and hydrophobic interactions between conserved alpha helices whose disruption by mutation damages the proteasome's ability to assemble. The assembly of the half-proteasomes, in turn, is initiated by the assembly of the α subunits into their heptameric ring, forming a template for the association of the corresponding pro-β ring. The assembly of α subunits has not been characterized. Only recently, the assembly process of the 19S regulatory particle has been elucidated to considerable extent. The 19S regulatory particle assembles as two distinct subcomponents, the base and the lid. Assembly of the base complex is facilitated by four assembly chaperones, Hsm3/S5b, Nas2/p27, Rpn14/PAAF1, and Nas6/gankyrin (names for yeast/mammals). These assembly chaperones bind to the AAA-ATPase subunits and their main function seems to be to ensure proper assembly of the heterohexameric AAA-ATPase ring. To date it is still under debate whether the base complex assembles separately, whether the assembly is templated by the 20S core particle, or whether alternative assembly pathways exist. In addition to the four assembly chaperones, the deubiquitinating enzyme Ubp6/Usp14 also promotes base assembly, but it is not essential. The lid assembles separately in a specific order and does not require assembly chaperones. Protein degradation process Ubiquitination and targeting Proteins are targeted for degradation by the proteasome with covalent modification of a lysine residue that requires the coordinated reactions of three enzymes. In the first step, a ubiquitin-activating enzyme (known as E1) hydrolyzes ATP and adenylylates a ubiquitin molecule. This is then transferred to E1's active-site cysteine residue in concert with the adenylylation of a second ubiquitin. This adenylylated ubiquitin is then transferred to a cysteine of a second enzyme, ubiquitin-conjugating enzyme (E2). In the last step, a member of a highly diverse class of enzymes known as ubiquitin ligases (E3) recognizes the specific protein to be ubiquitinated and catalyzes the transfer of ubiquitin from E2 to this target protein. A target protein must be labeled with at least four ubiquitin monomers (in the form of a polyubiquitin chain) before it is recognized by the proteasome lid. It is therefore the E3 that confers substrate specificity to this system. The number of E1, E2, and E3 proteins expressed depends on the organism and cell type, but there are many different E3 enzymes present in humans, indicating that there is a huge number of targets for the ubiquitin proteasome system. The mechanism by which a polyubiquitinated protein is targeted to the proteasome is not fully understood. A few high-resolution snapshots of the proteasome bound to a polyubiquitinated protein suggest that ubiquitin receptors might be coordinated with deubiquitinase Rpn11 for initial substrate targeting and engagement. Ubiquitin-receptor proteins have an N-terminal ubiquitin-like (UBL) domain and one or more ubiquitin-associated (UBA) domains. The UBL domains are recognized by the 19S proteasome caps and the UBA domains bind ubiquitin via three-helix bundles. These receptor proteins may escort polyubiquitinated proteins to the proteasome, though the specifics of this interaction and its regulation are unclear. The ubiquitin protein itself is 76 amino acids long and was named due to its ubiquitous nature, as it has a highly conserved sequence and is found in all known eukaryotic organisms. The genes encoding ubiquitin in eukaryotes are arranged in tandem repeats, possibly due to the heavy transcription demands on these genes to produce enough ubiquitin for the cell. It has been proposed that ubiquitin is the slowest-evolving protein identified to date. Ubiquitin contains seven lysine residues to which another ubiquitin can be ligated, resulting in different types of polyubiquitin chains. Chains in which each additional ubiquitin is linked to lysine 48 of the previous ubiquitin have a role in proteasome targeting, while other types of chains may be involved in other processes. Deubiquitylation Ubiquitin chains conjugated to a protein targeted for proteasomal degradation are normally removed by any one of the three proteasome-associated deubiquitylating enzymes (DUBs), which are Rpn11, Ubp6/USP14 and UCH37. This process recycles ubiquitin and is essential to maintain the ubiquitin reservoir in cells. Rpn11 is an intrinsic, stoichiometric subunit of the 19S regulatory particle and is essential for the function of 26S proteasome. The DUB activity of Rpn11 is enhanced in the proteasome as compared to its monomeric form. How Rpn11 removes a ubiquitin chain en bloc from a protein substrate was captured by an atomic structure of the substrate-engaged human proteasome in a conformation named EB. Interestingly, this structure also shows how the DUB activity is coupled to the substrate recognition by the proteasomal AAA-ATPase. In contrast to Rpn11, USP14 and UCH37 are the DUBs that do not always associated with the proteasome. In cells, about 10-40% of the proteasomes were found to have USP14 associated. Both Ubp6/USP14 and UCH37 are largely activated by the proteasome and exhibit a very low DUB activity alone. Once activated, USP14 was found to suppress proteasome function by its DUB activity and by inducing parallel pathways of proteasome conformational transitions, one of which turned out to directly prohibit substrate insertion into the AAA-ATPase, as intuitively observed by time-resolved cryogenic electron microscopy. It appears that USP14 regulates proteasome function at multiple checkpoints by both catalytically competing with Rpn11 and allosterically reprogramming the AAA-ATPase states, which is rather unexpected for a DUB. These observations imply that the proteasome regulation may depend on its dynamic transitions of conformational states. Unfolding and translocation After a protein has been ubiquitinated, it is recognized by the 19S regulatory particle in an ATP-dependent binding step. The substrate protein must then enter the interior of the 20S subunit to come in contact with the proteolytic active sites. Because the 20S particle's central channel is narrow and gated by the N-terminal tails of the α ring subunits, the substrates must be at least partially unfolded before they enter the core. The passage of the unfolded substrate into the core is called translocation and necessarily occurs after deubiquitination. However, the order in which substrates are deubiquitinated and unfolded is not yet clear. Which of these processes is the rate-limiting step in the overall proteolysis reaction depends on the specific substrate; for some proteins, the unfolding process is rate-limiting, while deubiquitination is the slowest step for other proteins. The extent to which substrates must be unfolded before translocation is suggested to be around 20 amino acid residues by the atomic structure of the substrate-engaged 26S proteasome in the deubiquitylation-compatible state, but substantial tertiary structure, and in particular nonlocal interactions such as disulfide bonds, are sufficient to inhibit degradation. The presence of intrinsically disordered protein segments of sufficient size, either at the protein terminus or internally, has also been proposed to facilitate efficient initiation of degradation. The gate formed by the α subunits prevents peptides longer than about four residues from entering the interior of the 20S particle. The ATP molecules bound before the initial recognition step are hydrolyzed before translocation. While energy is needed for substrate unfolding, it is not required for translocation. The assembled 26S proteasome can degrade unfolded proteins in the presence of a non-hydrolyzable ATP analog, but cannot degrade folded proteins, indicating that energy from ATP hydrolysis is used for substrate unfolding. Passage of the unfolded substrate through the opened gate occurs via facilitated diffusion if the 19S cap is in the ATP-bound state. The mechanism for unfolding of globular proteins is necessarily general, but somewhat dependent on the amino acid sequence. Long sequences of alternating glycine and alanine have been shown to inhibit substrate unfolding, decreasing the efficiency of proteasomal degradation; this results in the release of partially degraded byproducts, possibly due to the decoupling of the ATP hydrolysis and unfolding steps. Such glycine-alanine repeats are also found in nature, for example in silk fibroin; in particular, certain Epstein–Barr virus gene products bearing this sequence can stall the proteasome, helping the virus propagate by preventing antigen presentation on the major histocompatibility complex. Proteolysis The proteasome functions as an endoprotease. The mechanism of proteolysis by the β subunits of the 20S core particle is through a threonine-dependent nucleophilic attack. This mechanism may depend on an associated water molecule for deprotonation of the reactive threonine hydroxyl. Degradation occurs within the central chamber formed by the association of the two β rings and normally does not release partially degraded products, instead reducing the substrate to short polypeptides typically 7–9 residues long, though they can range from 4 to 25 residues, depending on the organism and substrate. The biochemical mechanism that determines product length is not fully characterized. Although the three catalytic β subunits have a common mechanism, they have slightly different substrate specificities, which are considered chymotrypsin-like, trypsin-like, and peptidyl-glutamyl peptide-hydrolyzing (PHGH)-like. These variations in specificity are the result of interatomic contacts with local residues near the active sites of each subunit. Each catalytic β subunit also possesses a conserved lysine residue required for proteolysis. Although the proteasome normally produces very short peptide fragments, in some cases these products are themselves biologically active and functional molecules. Certain transcription factors regulating the expression of specific genes, including one component of the mammalian complex NF-κB, are synthesized as inactive precursors whose ubiquitination and subsequent proteasomal degradation converts them to an active form. Such activity requires the proteasome to cleave the substrate protein internally, rather than processively degrading it from one terminus. It has been suggested that long loops on these proteins' surfaces serve as the proteasomal substrates and enter the central cavity, while the majority of the protein remains outside. Similar effects have been observed in yeast proteins; this mechanism of selective degradation is known as regulated ubiquitin/proteasome dependent processing (RUP). Ubiquitin-independent degradation Although most proteasomal substrates must be ubiquitinated before being degraded, there are some exceptions to this general rule, especially when the proteasome plays a normal role in the post-translational processing of the protein. The proteasomal activation of NF-κB by processing p105 into p50 via internal proteolysis is one major example. Some proteins that are hypothesized to be unstable due to intrinsically unstructured regions, are degraded in a ubiquitin-independent manner. The most well-known example of a ubiquitin-independent proteasome substrate is the enzyme ornithine decarboxylase. Ubiquitin-independent mechanisms targeting key cell cycle regulators such as p53 have also been reported, although p53 is also subject to ubiquitin-dependent degradation. Finally, structurally abnormal, misfolded, or highly oxidized proteins are also subject to ubiquitin-independent and 19S-independent degradation under conditions of cellular stress. Evolution The 20S proteasome is both ubiquitous and essential in eukaryotes and archaea. The bacterial order Actinomycetales, also share homologs of the 20S proteasome, whereas most bacteria possess heat shock genes hslV and hslU, whose gene products are a multimeric protease arranged in a two-layered ring and an ATPase. The hslV protein has been hypothesized to resemble the likely ancestor of the 20S proteasome. In general, HslV is not essential in bacteria, and not all bacteria possess it, whereas some protists possess both the 20S and the hslV systems. Many bacteria also possess other homologs of the proteasome and an associated ATPase, most notably ClpP and ClpX. This redundancy explains why the HslUV system is not essential. Sequence analysis suggests that the catalytic β subunits diverged earlier in evolution than the predominantly structural α subunits. In bacteria that express a 20S proteasome, the β subunits have high sequence identity to archaeal and eukaryotic β subunits, whereas the α sequence identity is much lower. The presence of 20S proteasomes in bacteria may result from lateral gene transfer, while the diversification of subunits among eukaryotes is ascribed to multiple gene duplication events. Cell cycle control Cell cycle progression is controlled by ordered action of cyclin-dependent kinases (CDKs), activated by specific cyclins that demarcate phases of the cell cycle. Mitotic cyclins, which persist in the cell for only a few minutes, have one of the shortest life spans of all intracellular proteins. After a CDK-cyclin complex has performed its function, the associated cyclin is polyubiquitinated and destroyed by the proteasome, which provides directionality for the cell cycle. In particular, exit from mitosis requires the proteasome-dependent dissociation of the regulatory component cyclin B from the mitosis promoting factor complex. In vertebrate cells, "slippage" through the mitotic checkpoint leading to premature M phase exit can occur despite the delay of this exit by the spindle checkpoint. Earlier cell cycle checkpoints such as post-restriction point check between G1 phase and S phase similarly involve proteasomal degradation of cyclin A, whose ubiquitination is promoted by the anaphase promoting complex (APC), an E3 ubiquitin ligase. The APC and the Skp1/Cul1/F-box protein complex (SCF complex) are the two key regulators of cyclin degradation and checkpoint control; the SCF itself is regulated by the APC via ubiquitination of the adaptor protein, Skp2, which prevents SCF activity before the G1-S transition. Individual components of the 19S particle have their own regulatory roles. Gankyrin, a recently identified oncoprotein, is one of the 19S subcomponents that also tightly binds the cyclin-dependent kinase CDK4 and plays a key role in recognizing ubiquitinated p53, via its affinity for the ubiquitin ligase MDM2. Gankyrin is anti-apoptotic and has been shown to be overexpressed in some tumor cell types such as hepatocellular carcinoma. Like eukaryotes, some archaea also use the proteasome to control cell cycle, specifically by controlling ESCRT-III-mediated cell division. Regulation of plant growth In plants, signaling by auxins, or phytohormones that order the direction and tropism of plant growth, induces the targeting of a class of transcription factor repressors known as Aux/IAA proteins for proteasomal degradation. These proteins are ubiquitinated by SCFTIR1, or SCF in complex with the auxin receptor TIR1. Degradation of Aux/IAA proteins derepresses transcription factors in the auxin-response factor (ARF) family and induces ARF-directed gene expression. The cellular consequences of ARF activation depend on the plant type and developmental stage, but are involved in directing growth in roots and leaf veins. The specific response to ARF derepression is thought to be mediated by specificity in the pairing of individual ARF and Aux/IAA proteins. Apoptosis Both internal and external signals can lead to the induction of apoptosis, or programmed cell death. The resulting deconstruction of cellular components is primarily carried out by specialized proteases known as caspases, but the proteasome also plays important and diverse roles in the apoptotic process. The involvement of the proteasome in this process is indicated by both the increase in protein ubiquitination, and of E1, E2, and E3 enzymes that is observed well in advance of apoptosis. During apoptosis, proteasomes localized to the nucleus have also been observed to translocate to outer membrane blebs characteristic of apoptosis. Proteasome inhibition has different effects on apoptosis induction in different cell types. In general, the proteasome is not required for apoptosis, although inhibiting it is pro-apoptotic in most cell types that have been studied. Apoptosis is mediated through disrupting the regulated degradation of pro-growth cell cycle proteins. However, some cell lines — in particular, primary cultures of quiescent and differentiated cells such as thymocytes and neurons — are prevented from undergoing apoptosis on exposure to proteasome inhibitors. The mechanism for this effect is not clear, but is hypothesized to be specific to cells in quiescent states, or to result from the differential activity of the pro-apoptotic kinase JNK. The ability of proteasome inhibitors to induce apoptosis in rapidly dividing cells has been exploited in several recently developed chemotherapy agents such as bortezomib and . Response to cellular stress In response to cellular stresses – such as infection, heat shock, or oxidative damage – heat shock proteins that identify misfolded or unfolded proteins and target them for proteasomal degradation are expressed. Both Hsp27 and Hsp90—chaperone proteins have been implicated in increasing the activity of the ubiquitin-proteasome system, though they are not direct participants in the process. Hsp70, on the other hand, binds exposed hydrophobic patches on the surface of misfolded proteins and recruits E3 ubiquitin ligases such as CHIP to tag the proteins for proteasomal degradation. The CHIP protein (carboxyl terminus of Hsp70-interacting protein) is itself regulated via inhibition of interactions between the E3 enzyme CHIP and its E2 binding partner. Similar mechanisms exist to promote the degradation of oxidatively damaged proteins via the proteasome system. In particular, proteasomes localized to the nucleus are regulated by PARP and actively degrade inappropriately oxidized histones. Oxidized proteins, which often form large amorphous aggregates in the cell, can be degraded directly by the 20S core particle without the 19S regulatory cap and do not require ATP hydrolysis or tagging with ubiquitin. However, high levels of oxidative damage increases the degree of cross-linking between protein fragments, rendering the aggregates resistant to proteolysis. Larger numbers and sizes of such highly oxidized aggregates are associated with aging. Dysregulation of the ubiquitin proteasome system may contribute to several neural diseases. It may lead to brain tumors such as astrocytomas. In some of the late-onset neurodegenerative diseases that share aggregation of misfolded proteins as a common feature, such as Parkinson's disease and Alzheimer's disease, large insoluble aggregates of misfolded proteins can form and then result in neurotoxicity, through mechanisms that are not yet well understood. Decreased proteasome activity has been suggested as a cause of aggregation and Lewy body formation in Parkinson's. This hypothesis is supported by the observation that yeast models of Parkinson's are more susceptible to toxicity from α-synuclein, the major protein component of Lewy bodies, under conditions of low proteasome activity. Impaired proteasomal activity may underlie cognitive disorders such as the autism spectrum disorders, and muscle and nerve diseases such as inclusion body myopathy. Role in the immune system The proteasome plays a straightforward but critical role in the function of the adaptive immune system. Peptide antigens are displayed by the major histocompatibility complex class I (MHC) proteins on the surface of antigen-presenting cells. These peptides are products of proteasomal degradation of proteins originated by the invading pathogen. Although constitutively expressed proteasomes can participate in this process, a specialized complex composed of proteins, whose expression is induced by interferon gamma, are the primary producers of peptides which are optimal in size and composition for MHC binding. These proteins whose expression increases during the immune response include the 11S regulatory particle, whose main known biological role is regulating the production of MHC ligands, and specialized β subunits called β1i, β2i, and β5i with altered substrate specificity. The complex formed with the specialized β subunits is known as the immunoproteasome. Another β5i variant subunit, β5t, is expressed in the thymus, leading to a thymus-specific "thymoproteasome" whose function is as yet unclear. The strength of MHC class I ligand binding is dependent on the composition of the ligand C-terminus, as peptides bind by hydrogen bonding and by close contacts with a region called the "B pocket" on the MHC surface. Many MHC class I alleles prefer hydrophobic C-terminal residues, and the immunoproteasome complex is more likely to generate hydrophobic C-termini. Due to its role in generating the activated form of NF-κB, an anti-apoptotic and pro-inflammatory regulator of cytokine expression, proteasomal activity has been linked to inflammatory and autoimmune diseases. Increased levels of proteasome activity correlate with disease activity and have been implicated in autoimmune diseases including systemic lupus erythematosus and rheumatoid arthritis. The proteasome is also involved in Intracellular antibody-mediated proteolysis of antibody-bound virions. In this neutralisation pathway, TRIM21 (a protein of the tripartite motif family) binds with immunoglobulin G to direct the virion to the proteasome where it is degraded. Proteasome inhibitors Proteasome inhibitors have effective anti-tumor activity in cell culture, inducing apoptosis by disrupting the regulated degradation of pro-growth cell cycle proteins. This approach of selectively inducing apoptosis in tumor cells has proven effective in animal models and human trials. Lactacystin, a natural product synthesized by Streptomyces bacteria, was the first non-peptidic proteasome inhibitor discovered and is widely used as a research tool in biochemistry and cell biology. Lactacystin was licensed to Myogenics/Proscript, which was acquired by Millennium Pharmaceuticals, now part of Takeda Pharmaceuticals. Lactacystin covalently modifies the amino-terminal threonine of catalytic β subunits of the proteasome, particularly the β5 subunit responsible for the proteasome's chymotrypsin-like activity. This discovery helped to establish the proteasome as a mechanistically novel class of protease: an amino-terminal threonine protease. Bortezomib (Boronated MG132), a molecule developed by Millennium Pharmaceuticals and marketed as Velcade, is the first proteasome inhibitor to reach clinical use as a chemotherapy agent. Bortezomib is used in the treatment of multiple myeloma. Notably, multiple myeloma has been observed to result in increased proteasome-derived peptide levels in blood serum that decrease to normal levels in response to successful chemotherapy. Studies in animals have indicated that bortezomib may also have clinically significant effects in pancreatic cancer. Preclinical and early clinical studies have been started to examine bortezomib's effectiveness in treating other B-cell-related cancers, particularly some types of non-Hodgkin's lymphoma. Clinical results also seem to justify use of proteasome inhibitor combined with chemotherapy, for B-cell acute lymphoblastic leukemia Proteasome inhibitors can kill some types of cultured leukemia cells that are resistant to glucocorticoids. The molecule ritonavir, marketed as Norvir, was developed as a protease inhibitor and used to target HIV infection. However, it has been shown to inhibit proteasomes as well as free proteases; to be specific, the chymotrypsin-like activity of the proteasome is inhibited by ritonavir, while the trypsin-like activity is somewhat enhanced. Studies in animal models suggest that ritonavir may have inhibitory effects on the growth of glioma cells. Proteasome inhibitors have also shown promise in treating autoimmune diseases in animal models. For example, studies in mice bearing human skin grafts found a reduction in the size of lesions from psoriasis after treatment with a proteasome inhibitor. Inhibitors also show positive effects in rodent models of asthma. Labeling and inhibition of the proteasome is also of interest in laboratory settings for both in vitro and in vivo study of proteasomal activity in cells. The most commonly used laboratory inhibitors are lactacystin and the peptide aldehyde MG132 initially developed by Goldberg lab. Fluorescent inhibitors have also been developed to specifically label the active sites of the assembled proteasome. Clinical significance The proteasome and its subunits are of clinical significance for at least two reasons: (1) a compromised complex assembly or a dysfunctional proteasome can be associated with the underlying pathophysiology of specific diseases, and (2) they can be exploited as drug targets for therapeutic interventions. More recently, more effort has been made to consider the proteasome for the development of novel diagnostic markers and strategies. An improved and comprehensive understanding of the pathophysiology of the proteasome should lead to clinical applications in the future. The proteasomes form a pivotal component for the ubiquitin–proteasome system (UPS) and corresponding cellular Protein Quality Control (PQC). Protein ubiquitination and subsequent proteolysis and degradation by the proteasome are important mechanisms in the regulation of the cell cycle, cell growth and differentiation, gene transcription, signal transduction and apoptosis. Proteasome defects lead to reduced proteolytic activity and the accumulation of damaged or misfolded proteins, which may contribute to neurodegenerative disease, cardiovascular diseases, inflammatory responses and autoimmune diseases, and systemic DNA damage responses leading to malignancies. Research has implicated UPS defects in the pathogenesis of neurodegenerative and myodegenerative disorders, including Alzheimer's disease, Parkinson's disease and Pick's disease, amyotrophic lateral sclerosis (ALS), Huntington's disease, Creutzfeldt–Jakob disease, and motor neuron diseases, polyglutamine (PolyQ) diseases, muscular dystrophies and several rare forms of neurodegenerative diseases associated with dementia. As part of the ubiquitin–proteasome system (UPS), the proteasome maintains cardiac protein homeostasis and thus plays a significant role in cardiac ischemic injury, ventricular hypertrophy and heart failure. Additionally, evidence is accumulating that the UPS plays an essential role in malignant transformation. UPS proteolysis plays a major role in responses of cancer cells to stimulatory signals that are critical for the development of cancer. Accordingly, gene expression by degradation of transcription factors, such as p53, c-jun, c-Fos, NF-κB, c-Myc, HIF-1α, MATα2, STAT3, sterol-regulated element-binding proteins and androgen receptors are all controlled by the UPS and thus involved in the development of various malignancies. Moreover, the UPS regulates the degradation of tumor suppressor gene products such as adenomatous polyposis coli (APC) in colorectal cancer, retinoblastoma (Rb). and von Hippel–Lindau tumor suppressor (VHL), as well as a number of proto-oncogenes (Raf, Myc, Myb, Rel, Src, Mos, ABL). The UPS is also involved in the regulation of inflammatory responses. This activity is usually attributed to the role of proteasomes in the activation of NF-κB which further regulates the expression of pro inflammatory cytokines such as TNF-α, IL-β, IL-8, adhesion molecules (ICAM-1, VCAM-1, P-selectin) and prostaglandins and nitric oxide (NO). Additionally, the UPS also plays a role in inflammatory responses as regulators of leukocyte proliferation, mainly through proteolysis of cyclines and the degradation of CDK inhibitors. Lastly, autoimmune disease patients with SLE, Sjögren syndrome and rheumatoid arthritis (RA) predominantly exhibit circulating proteasomes which can be applied as clinical biomarkers.	Biology and health sciences	Organelles	Biology
24619	https://en.wikipedia.org/wiki/Pottery	Pottery	Pottery is the process and the products of forming vessels and other objects with clay and other raw materials, which are fired at high temperatures to give them a hard and durable form. The place where such wares are made by a potter is also called a pottery (plural potteries). The definition of pottery, used by the ASTM International, is "all fired ceramic wares that contain clay when formed, except technical, structural, and refractory products". End applications include tableware, decorative ware, sanitary ware, and in technology and industry such as electrical insulators and laboratory ware. In art history and archaeology, especially of ancient and prehistoric periods, pottery often means only vessels, and sculpted figurines of the same material are called terracottas. Pottery is one of the oldest human inventions, originating before the Neolithic period, with ceramic objects such as the Gravettian culture Venus of Dolní Věstonice figurine discovered in the Czech Republic dating back to 29,000–25,000 BC. However, the earliest known pottery vessels were discovered in Jiangxi, China, which date back to 18,000 BC. Other early Neolithic and pre-Neolithic pottery artifacts have been found, in Jōmon Japan (10,500 BC), the Russian Far East (14,000 BC), Sub-Saharan Africa (9,400 BC), South America (9,000s–7,000s BC), and the Middle East (7,000s–6,000s BC). Pottery is made by forming a clay body into objects of a desired shape and heating them to high temperatures (600–1600 °C) in a bonfire, pit or kiln, which induces reactions that lead to permanent changes including increasing the strength and rigidity of the object. Much pottery is purely utilitarian, but some can also be regarded as ceramic art. An article can be decorated before or after firing. Pottery is traditionally divided into three types: earthenware, stoneware and porcelain. All three may be glazed and unglazed. All may also be decorated by various techniques. In many examples the group a piece belongs to is immediately visually apparent, but this is not always the case; for example fritware uses no or little clay, so falls outside these groups. Historic pottery of all these types is often grouped as either "fine" wares, relatively expensive and well-made, and following the aesthetic taste of the culture concerned, or alternatively "coarse", "popular", "folk" or "village" wares, mostly undecorated, or simply so, and often less well-made. Cooking in pottery became less popular once metal pots became available, but is still used for dishes that benefit from the qualities of pottery cooking, typically slow cooking in an oven, such as biryani, cassoulet, daube, tagine, jollof rice, kedjenou, cazuela and types of baked beans. Main types Earthenware The earliest forms of pottery were made from clays that were fired at low temperatures, initially in pit-fires or in open bonfires. They were hand formed and undecorated. Earthenware can be fired as low as 600 °C, and is normally fired below 1200 °C. Because unglazed earthenware is porous, it has limited utility for the storage of liquids or as tableware. However, earthenware has had a continuous history from the Neolithic period to today. It can be made from a wide variety of clays, some of which fire to a buff, brown or black colour, with iron in the constituent minerals resulting in a reddish-brown. Reddish coloured varieties are called terracotta, especially when unglazed or used for sculpture. The development of ceramic glaze made impermeable pottery possible, improving the popularity and practicality of pottery vessels. Decoration has evolved and developed through history. Stoneware Stoneware is pottery that has been fired in a kiln at a relatively high temperature, from about 1,100 °C to 1,200 °C, and is stronger and non-porous to liquids. The Chinese, who developed stoneware very early on, classify this together with porcelain as high-fired wares. In contrast, stoneware could only be produced in Europe from the late Middle Ages, as European kilns were less efficient, and the right type of clay less common. It remained a speciality of Germany until the Renaissance. Stoneware is very tough and practical, and much of it has always been utilitarian, for the kitchen or storage rather than the table. But "fine" stoneware has been important in China, Japan and the West, and continues to be made. Many utilitarian types have also come to be appreciated as art. Porcelain Porcelain is made by heating materials, generally including kaolin, in a kiln to temperatures between . This is higher than used for the other types, and achieving these temperatures was a long struggle, as well as realizing what materials were needed. The toughness, strength and translucence of porcelain, relative to other types of pottery, arises mainly from vitrification and the formation of the mineral mullite within the body at these high temperatures. Although porcelain was first made in China, the Chinese traditionally do not recognise it as a distinct category, grouping it with stoneware as "high-fired" ware, opposed to "low-fired" earthenware. This confuses the issue of when it was first made. A degree of translucency and whiteness was achieved by the Tang dynasty (AD 618–906), and considerable quantities were being exported. The modern level of whiteness was not reached until much later, in the 14th century. Porcelain was also made in Korea and in Japan from the end of the 16th century, after suitable kaolin was located in those countries. It was not made effectively outside East Asia until the 18th century. Archaeology The study of pottery can help to provide an insight into past cultures. Fabric analysis (see section below), used to analyse the fabric of pottery, is important part of archaeology for understanding the archaeological culture of the excavated site by studying the fabric of artifacts, such as their usage, source material composition, decorative pattern, color of patterns, etc. This helps to understand characteristics, sophistication, habits, technology, tools, trade, etc. of the people who made and used the pottery. Carbon dating reveals the age. Sites with similar pottery characteristics have the same culture, those sites which have distinct cultural characteristics but with some overlap are indicative of cultural exchange such as trade or living in vicinity or continuity of habitation, etc. Examples are black and red ware, redware, Sothi-Siswal culture and Painted Grey Ware culture. The six fabrics of Kalibangan is a good example of use of fabric analysis in identifying a differentiated culture which was earlier thought to be typical Indus Valley civilisation (IVC) culture. Pottery is durable, and fragments, at least, often survive long after artifacts made from less-durable materials have decayed past recognition. Combined with other evidence, the study of pottery artefacts is helpful in the development of theories on the organisation, economic condition and the cultural development of the societies that produced or acquired pottery. The study of pottery may also allow inferences to be drawn about a culture's daily life, religion, social relationships, attitudes towards neighbours, attitudes to their own world and even the way the culture understood the universe. It is valuable to look into pottery as an archaeological record of potential interaction between peoples. When pottery is placed within the context of linguistic and migratory patterns, it becomes an even more prevalent category of social artifact. As proposed by Olivier P. Gosselain, it is possible to understand ranges of cross-cultural interaction by looking closely at the chaîne opératoire of ceramic production. The methods used to produce pottery in early Sub-Saharan Africa are divisible into three categories: techniques visible to the eye (decoration, firing and post-firing techniques), techniques related to the materials (selection or processing of clay, etc.), and techniques of molding or fashioning the clay. These three categories can be used to consider the implications of the reoccurrence of a particular sort of pottery in different areas. Generally, the techniques that are easily visible (the first category of those mentioned above) are thus readily imitated, and may indicate a more distant connection between groups, such as trade in the same market or even relatively close settlements. Techniques that require more studied replication (i.e., the selection of clay and the fashioning of clay) may indicate a closer connection between peoples, as these methods are usually only transmissible between potters and those otherwise directly involved in production. Such a relationship requires the ability of the involved parties to communicate effectively, implying pre-existing norms of contact or a shared language between the two. Thus, the patterns of technical diffusion in pot-making that are visible via archaeological findings also reveal patterns in societal interaction. Chronologies based on pottery are often essential for dating non-literate cultures and are often of help in the dating of historic cultures as well. Trace-element analysis, mostly by neutron activation, allows the sources of clay to be accurately identified and the thermoluminescence test can be used to provide an estimate of the date of last firing. Examining sherds from prehistory, scientists learned that during high-temperature firing, iron materials in clay record the state of the Earth's magnetic field at that moment. Fabric analysis The "clay body" is also called the "paste" or the "fabric", which consists of 2 things, the "clay matrix" – composed of grains of less than 0.02 mm grains which can be seen using the high-powered microscopes or a scanning electron microscope, and the "clay inclusions" – which are larger grains of clay and could be seen with the naked eye or a low-power binocular microscope. For geologists, fabric analysis means spatial arrangement of minerals in a rock. For Archaeologists, the "fabric analysis" of pottery entails the study of clay matrix and inclusions in the clay body as well as the firing temperature and conditions. Analysis is done to examine the following 3 in detail: How pottery was made e.g. material, design such as shape and style, etc. Its decorations, such as patterns, colors of patterns, slipped (glazing) or unslipped decoration Evidence of type of use. The Six fabrics of Kalibangan is a good example of fabric analysis. Clay bodies and raw materials Body, or clay body, is the material used to form pottery. Thus a potter might prepare, or order from a supplier, such an amount of earthenware body, stoneware body or porcelain body. The compositions of clay bodies varies considerably, and include both prepared and 'as dug'; the former being by far the dominant type for studio and industry. The properties also vary considerably, and include plasticity and mechanical strength before firing; the firing temperature needed to mature them; properties after firing, such as permeability, mechanical strength and colour. There can be regional variations in the properties of raw materials used for pottery, and these can lead to wares that are unique in character to a locality. The main ingredient of the body is clay. Some different types used for pottery include: Kaolin, sometimes referred to as china clay, is a key ingredient in porcelain, which was first used in China around the 7th and 8th centuries. Ball clay: An extremely plastic, fine grained sedimentary clay, which may contain some organic matter. Fire clay: A clay having a slightly lower percentage of fluxes than kaolin, but usually quite plastic. It is highly heat resistant form of clay which can be combined with other clays to increase the firing temperature and may be used as an ingredient to make stoneware type bodies. Stoneware clay: Suitable for creating stoneware. Has many of the characteristics between fire clay and ball clay, having finer grain, like ball clay but is more heat resistant like fire clays. Common red clay and shale clay have vegetable and ferric oxide impurities which make them useful for bricks, but are generally unsatisfactory for pottery except under special conditions of a particular deposit. Bentonite: An extremely plastic clay which can be added in small quantities to short clay to increase the plasticity. It is common for clays and other raw materials to be mixed to produce clay bodies suited to specific purposes. Various mineral processing techniques are often utilised before mixing the raw materials, with comminution being effectively universal for non-clay materials. Examples of non-clay materials include: Feldspar, act as fluxes which lower the vitrification temperature of bodies. Quartz, an important role is to attenuate drying shrinkage. Nepheline syenite, an alternative to feldspar. Calcined alumina, can enhance the fired properties of a body. Chamotte, also called grog, is fired clay which it is crushed, and sometimes then milled. Helps attenuate drying shrinkage. Bone ash, produced by the calcination of animal bone. A key raw material for bone china. Frit, produced made by quenching and breaking up a glass of a specific composition. Can be used at low additions in some bodies, but common uses include as components of a glaze or enamel, or for the body of fritware, when it usually mixed with larger quantities of quartz sand. Various others at low levels of addition such as dolomite, limestone, talc and wollastonite. Production The production of pottery includes the following stages: Preparing the clay body. Shaping Drying Firing Glazing and decorating. (this can be undertaken prior to firing. Also, additional firing stages after decoration may be needed.) Shaping Before being shaped, clay must be prepared. This may include kneading to ensure an even moisture content throughout the body. Air trapped within the clay body needs to be removed, or de-aired, and can be accomplished either by a machine called a vacuum pug or manually by wedging. Wedging can also help produce an even moisture content. Once a clay body has been kneaded and de-aired or wedged, it is shaped by a variety of techniques, which include: Hand-building: This is the earliest forming method. Wares can be constructed by hand from coils of clay, combining flat slabs of clay, or pinching solid balls of clay or some combination of these. Parts of hand-built vessels are often joined with the aid of slip. Some studio potters find hand-building more conducive for one-of-a-kind works of art. The potter's wheel: In a process called "throwing" (coming from the Old English word thrawan which means to twist or turn,) a ball of clay is placed in the centre of a turntable, called the wheel-head, which the potter rotates with a stick, with foot power or with a variable-speed electric motor. During the process of throwing, the wheel rotates while the solid ball of soft clay is pressed, squeezed and pulled gently upwards and outwards into a hollow shape. Skill and experience are required to throw pots of an acceptable standard and, while the ware may have high artistic merit, the reproducibility of the method is poor. Because of its inherent limitations, throwing can only be used to create wares with radial symmetry on a vertical axis. Press moulding: a simple technique of shaping by manually pressing a lump of clay body into a porous mould. Granulate pressing: a highly automated technique of shaping by pressing clay body in a semi-dry and granulated form in a mould. The body is pressed into the mould by a porous die through which water is pumped at high pressure. The fine, free flowing granulated body is prepared by spray drying a high-solids content slip. Granulate pressing, also known as dust pressing, is widely used in the manufacture of ceramic tiles and, increasingly, of plates. Jiggering and jolleying: These operations are carried out on the potter's wheel and allow the time taken to bring wares to a standardized form to be reduced. Jiggering is the operation of bringing a shaped tool into contact with the plastic clay of a piece under construction, the piece itself being set on a rotating plaster mould on the wheel. The jigger tool shapes one face while the mould shapes the other. Jiggering is used only in the production of flat wares, such as plates, but a similar operation, jolleying, is used in the production of hollow-wares such as cups. Jiggering and jolleying have been used in the production of pottery since at least the 18th century. In large-scale factory production, jiggering and jolleying are usually automated, which allows the operations to be carried out by semi-skilled labour. Roller-head machine: This machine is for shaping wares on a rotating mould, as in jiggering and jolleying, but with a rotary shaping tool replacing the fixed profile. The rotary shaping tool is a shallow cone having the same diameter as the ware being formed and shaped to the desired form of the back of the article being made. Wares may in this way be shaped, using relatively unskilled labour, in one operation at a rate of about twelve pieces per minute, though this varies with the size of the articles being produced. Developed in the UK just after World War II by the company Service Engineers, roller-heads were quickly adopted by manufacturers around the world; it remains the dominant method for producing both flatware and holloware, such as plates and mugs. Pressure casting: Is a development of traditional slipcasting. Specially developed polymeric materials allow a mould to be subject to application external pressures of up to 4.0 MPa – so much higher than slip casting in plaster moulds where the capillary forces correspond to a pressure of around 0.1–0.2 MPa. The high pressure leads to much faster casting rates and, hence, faster production cycles. Furthermore, the application of high pressure air through the polymeric moulds upon demoulding the cast means a new casting cycle can be started immediately in the same mould, unlike plaster moulds which require lengthy drying times. The polymeric materials have much greater durability than plaster and, therefore, it is possible to achieve shaped products with better dimensional tolerances and much longer mould life. Pressure casting was developed in the 1970s for the production of sanitaryware although, more recently, it has been applied to tableware. RAM pressing: This is used to shape ware by pressing a bat of prepared clay body into a required shape between two porous moulding plates. After pressing, compressed air is blown through the porous mould plates to release the shaped wares. Slip casting: This is suited to the making of shapes that cannot be formed by other methods. A liquid slip, made by mixing clay body with water, is poured into a highly absorbent plaster mould. Water from the slip is absorbed into the mould leaving a layer of clay body covering its internal surfaces and taking its internal shape. Excess slip is poured out of the mould, which is then split open and the moulded object removed. Slip casting is widely used in the production of sanitaryware and is also used for making other complex shaped ware such as teapots and figurines. Injection moulding: This is a shape-forming process adapted for the tableware industry from the method long established for the forming of thermoplastic and some metal components. It has been called Porcelain Injection Moulding, or PIM. Suited to the mass production of complex-shaped articles, one significant advantage of the technique is that it allows the production of a cup, including the handle, in a single process, and thereby eliminates the handle-fixing operation and produces a stronger bond between cup and handle. The feed to the mould die is a mix of approximately 50 to 60 per cent unfired body in powder form, together with 40 to 50 per cent organic additives composed of binders, lubricants and plasticisers. The technique is not as widely used as other shaping methods. 3D printing: There are two methods. One involves the layered deposition of soft clay body similar to fused deposition modelling (FDM), and the other uses powder binding techniques where clay body in dry powder form is fused together layer upon layer with a liquid. Injection moulding of ceramic tableware has been developed, though it has yet to be fully commercialised. Drying Prior to firing, the water in an article needs to be removed. A number of different stages, or conditions of the article, can be identified: Greenware refers to unfired objects at any stage of dryness, but is most often used to refer to objects ready to be fired. At sufficient moisture content, bodies at this stage are in their most plastic form (as they are soft and malleable, and hence can be easily deformed by handling). Prior to firing, any state of clay may be hydrated or dehydrated into any other unfired stage. Plastic, also known as wet, refers to clay that is malleable and sufficiently wet to shape by hand or on a potter's wheel, but strong enough to hold its shape. At this stage the clay has between 20% and 25% moisture content. This is the stage most commercial clays are sold at, and at which most of the shaping process is done. Leather-hard refers to a clay body that has been dried partially. At this stage the clay object has approximately 15% moisture content. Clay bodies at this stage are very firm and only slightly pliable. Trimming and handle attachment often occurs at the leather-hard state. Bone-dry refers to clay bodies when they reach a moisture content at or near 0%. At that moisture content, the item is ready to be fired. Additionally, the piece is extremely brittle at this stage and must be handled with care. Firing Firing produces permanent and irreversible chemical and physical changes in the body. It is only after firing that the article or material is pottery. In lower-fired pottery, the changes include sintering, the fusing together of coarser particles in the body at their points of contact with each other. In the case of porcelain, where higher firing-temperatures are used, the physical, chemical and mineralogical properties of the constituents in the body are greatly altered. In all cases, the reason for firing is to permanently harden the wares, and the firing regime must be appropriate to the materials used. Temperature As a rough guide, modern earthenwares are normally fired at temperatures in the range of about to ; stonewares at between about to ; and porcelains at between about to . Historically, reaching high temperatures was a long-lasting challenge, and earthenware can be fired effectively as low as , achievable in primitive pit firing. The time spent at any particular temperature is also important, the combination of heat and time is known as heatwork. Kilns can be monitored by pyrometers, thermocouples and pyrometric devices. Atmosphere The atmosphere within a kiln during firing can affect the appearance of the body and glaze. Key to this is the differing colours of the various oxides of iron, such as iron(III) oxide (also known as ferric oxide or Fe2O3) which is associated with brown-red colours, whilst iron(II) oxide (also known as ferrous oxide or FeO) is associated with much darker colours, including black. The oxygen concentration in the kiln influences the type, and relative proportions, of these iron oxides in fired the body and glaze: for example, where there is a lack of oxygen during firing the associated carbon monoxide (CO) will readily react with oxygen in Fe2O3 in the raw materials and cause it to be reduced to FeO. An oxygen deficient condition, called a reducing atmosphere, is generated by preventing the complete combustion of the kiln fuel; this is achieved by deliberately restricting the supply of air or by supplying an excess of fuel. Methods Firing pottery can be done using a variety of methods, with a kiln being the usual firing method. Both the maximum temperature and the duration of firing influences the final characteristics of the ceramic. Thus, the maximum temperature within a kiln is often held constant for a period of time to soak the wares to produce the maturity required in the body of the wares. Kilns may be heated by burning combustible materials, such as wood, coal and gas, or by electricity. The use of microwave energy has been investigated. When used as fuels, coal and wood can introduce smoke, soot and ash into the kiln which can affect the appearance of unprotected wares. For this reason, wares fired in wood- or coal-fired kilns are often placed in the kiln in saggars, ceramic boxes, to protect them. Modern kilns fuelled by gas or electricity are cleaner and more easily controlled than older wood- or coal-fired kilns and often allow shorter firing times to be used. Niche techniques include: In a Western adaptation of traditional Japanese raku ware firing, wares are removed from the kiln while hot and smothered in ashes, paper or woodchips which produces a distinctive carbonised appearance. This technique is also used in Malaysia in creating traditional labu sayung. In Mali, a firing mound is used rather than a brick or stone kiln. Unfired pots are first brought to the place where a mound will be built, customarily by the women and girls of the village. The mound's foundation is made by placing sticks on the ground, then: Stages Biscuit (or bisque) refers to the clay after the object is shaped to the desired form and fired in the kiln for the first time, known as "bisque fired" or "biscuit fired". This firing results in both chemical and physical changes to the minerals of the clay body. Glaze fired is the final stage of some pottery making, or glost fired. A glaze may be applied to the biscuit ware and the object can be decorated in several ways. After this the object is "glazed fired", which causes the glaze material to melt, then adhere to the object. Depending on the temperature schedule the glaze firing may also further mature the body as chemical and physical changes continue. Decorating Pottery may be decorated in many different ways. Some decoration can be done before or after the firing, and may be undertaken before or after glazing. Methods Painting has been used since early prehistoric times, and can be very elaborate. The painting is often applied to pottery that has been fired once, and may then be overlaid with a glaze afterwards. Many pigments change colour when fired, and the painter must allow for this. Glaze: Perhaps the most common form of decoration, that also serves as protection to the pottery, by being tougher and keeping liquid from penetrating the pottery. Glaze may be colourless, especially over painting, or coloured and opaque. Crystalline glaze: acharacterised by crystalline clusters of various shapes and colours embedded in a more uniform and opaque glaze. Produced by the slow cooling of the glost fire. Carving: Pottery vessels may be decorated by shallow carving of the clay body, typically with a knife or similar instrument used on the wheel. This is common in Chinese porcelain of the classic periods. Burnishing: The surface of pottery wares may be burnished prior to firing by rubbing with a suitable instrument of wood, steel or stone to produce a polished finish that survives firing. It is possible to produce very highly polished wares when fine clays are used or when the polishing is carried out on wares that have been partially dried and contain little water, though wares in this condition are extremely fragile and the risk of breakage is high. Terra Sigillata is an ancient form of decorating ceramics that was first developed in Ancient Greece. Lithography, also called litho, although the alternative names of transfer print or "decal" are also common. These are used to apply designs to articles. The litho comprises three layers: the colour, or image, layer which comprises the decorative design; the cover coat, a clear protective layer, which may incorporate a low-melting glass; and the backing paper on which the design is printed by screen printing or lithography. There are various methods of transferring the design while removing the backing-paper, some of which are suited to machine application. Banding is the application by hand or by machine of a band of colour to the edge of a plate or cup. Also known as "lining", this operation is often carried out on a potter's wheel. Agateware: named after its resemblance to the mineral agate, is produced by partially blending clays of differing colours. In Japan the term "neriage" is used, whilst in China, where such things have been made since at least the Tang dynasty, they are called "marbled" wares. Engobe: a clay slip is used to coat the surface of pottery, usually before firing. Its purpose is often decorative though it can also be used to mask undesirable features in the clay to which it is applied. The engobe may be applied by painting or by dipping to provide a uniform, smooth, coating. Such decoration is characteristic of slipware. For sgraffito decoration a layer of engobe is scratched through to reveal the underlying clay. Gold: Decoration with gold is used on some high quality ware. Different methods exist for its application, including: Best gold – a suspension of gold powder in essential oils mixed with a flux and a mercury salt extended. This can be applied by a painting technique. From the kiln, the decoration is dull and requires burnishing to reveal the full colour Acid Gold – a form of gold decoration developed in the early 1860s at the English factory of Mintons Ltd. The glazed surface is etched with diluted hydrofluoric acid prior to application of the gold. The process demands great skill and is used for the decoration only of ware of the highest class. Bright Gold – consists of a solution of gold sulphoresinate together with other metal resonates and a flux. The name derives from the appearance of the decoration immediately after removal from the kiln as it requires no burnishing Mussel Gold – an old method of gold decoration. It was made by rubbing together gold leaf, sugar and salt, followed by washing to remove solubles Underglaze decoration is applied, by a number of techniques, onto ware before it is glazed, an example is blue and white wares. Can be applied by a number of techniques. In-glaze decoration, is applied on the surface of the glaze before the glost firing. On-glaze decoration is applied on top of the already fired, glazed surface, and then fixed in a second firing at a relatively low temperature. Glazing Glaze is a glassy coating on pottery, and reasons to use it include decoration, ensuring the item is impermeable to liquids, and minimizing the adherence of pollutants. Glaze may be applied by spraying, dipping, trailing or brushing on an aqueous suspension of the unfired glaze. The colour of a glaze after it has been fired may be significantly different from before firing. To prevent glazed wares sticking to kiln furniture during firing, either a small part of the object being fired (for example, the foot) is left unglazed or, alternatively, special refractory "spurs" are used as supports. These are removed and discarded after the firing. Some specialised glazing techniques include: Salt-glazing – common salt is introduced to the kiln during the firing process. The high temperatures cause the salt to volatilise, depositing it on the surface of the ware to react with the body to form a sodium aluminosilicate glaze. In the 17th and 18th centuries, salt-glazing was used in the manufacture of domestic pottery. Now, except for use by some studio potters, the process is obsolete. The last large-scale application before its demise in the face of environmental clean air restrictions was in the production of salt-glazed sewer-pipes. Ash glazing – ash from the combustion of plant matter has been used as the flux component of glazes. The source of the ash was generally the combustion waste from the fuelling of kilns although the potential of ash derived from arable crop wastes has been investigated. Ash glazes are of historical interest in the Far East although there are reports of small-scale use in other locations such as the Catawba Valley Pottery in the United States. They are now limited to small numbers of studio potters who value the unpredictability arising from the variable nature of the raw material. Types of Glazing in Pottery. Glazing in pottery is the process of applying a coating or layer of material to ceramics that, when fired, forms a vitreous or glass-like surface. Glazes enhance the aesthetic appeal of pottery, provide a waterproof barrier, and improve its durability. Below are the major types of glazing commonly used in pottery: 1. Glossy Glaze - Produces a shiny, reflective surface. - Highlights intricate patterns and textures. - Often used for decorative purposes. 2. Matte Glaze - Provides a smooth, non-reflective finish. - Suitable for modern and minimalist designs. - Ideal for functional wares like plates and mugs, as it minimizes glare. 3. Transparent Glaze - Can be glossy or matte. - Allows the underlying decoration or texture of the pottery to show through. - Often used over underglaze decorations. 4. Opaque Glaze - Completely covers the surface of the pottery, hiding any underlying texture or decoration. - Useful for creating uniform, bold colors. 5. Celadon Glaze - A translucent glaze, usually in shades of green or blue. - Originated in China and is popular in East Asian ceramics. - Accentuates carved or textured designs beneath the glaze. 6. Ash Glaze - Made from natural wood ash. - Creates unpredictable and unique textures during firing. - Often used in wood-fired kilns to achieve traditional aesthetics. 7. Salt Glaze - Achieved by introducing salt into a hot kiln during firing. - Produces a textured, orange-peel-like surface. - Common in traditional stoneware pottery. 8. Shino Glaze - A traditional Japanese glaze made from feldspar and clay. - Produces warm earthy tones, such as orange, red, and brown. - Works well in wood-fired pottery. 9. Satin Glaze - Offers a finish between matte and glossy. - Combines elegance with functionality. - Popular for tableware and decorative pieces. 10. Textured Glaze - Designed to create raised or recessed textures. - Adds visual and tactile interest to pottery. 11. Raku Glaze - Specifically designed for raku firing, where pottery is removed from the kiln while still hot. - Produces crackled textures and unique, unpredictable patterns. - Often used for artistic or decorative purposes. 12. Metallic or Lustre Glaze - Contains metal oxides or compounds, producing a metallic or iridescent finish. - Adds a luxurious appearance to pottery. - Requires specialized firing techniques. 13. Lead Glaze - Historically used for its shiny and smooth finish. - Toxicity concerns have reduced its use in modern pottery. 14. Feldspathic Glaze - Made from feldspar minerals. - Creates a glossy, natural finish. - Common in high-fired ceramics. 15. Engobe or Slip Glaze - A combination of slip (liquid clay) and glaze. - Used to decorate or add texture before applying a final glaze layer. Conclusion Glazes are a vital part of pottery, offering endless possibilities for functional and artistic expression. Each type of glaze interacts uniquely with the clay body and firing conditions, allowing potters to create diverse effects and finishes. Health and environmental issues Although many of the environmental effects of pottery production have existed for millennia, some of these have been amplified with modern technology and scales of production. The principal factors for consideration fall into two categories: Effects on workers: Notable risks include silicosis, heavy metal poisoning, poor indoor air quality, dangerous sound levels and possible over-illumination. Effects on the general environment. Historically, lead poisoning (plumbism) was a significant health concern to those glazing pottery. This was recognised at least as early as the nineteenth century. The first legislation in the UK to limit pottery workers exposure to lead was included in the Factories Act Extension Act in 1864, with further introduced in 1899. Silicosis is an occupational lung disease caused by inhaling large amounts of crystalline silica dust, usually over many years. Workers in the ceramic industry can develop it due to exposure to silica dust in the raw materials; colloquially it has been known as 'Potter's rot'. Less than 10 years after its introduction, in 1720, as a raw material to the British ceramics industry the negative effects of calcined flint on the lungs of workers had been noted. In one study reported in 2022, of 106 UK pottery workers 55 per cent had at least some stage of silicosis. Exposure to siliceous dusts is reduced by either processing and using the source materials as aqueous suspension or as damp solids, or by the use of dust control measures such as local exhaust ventilation. These have been mandated by legislation, such as The Pottery (Health and Welfare) Special Regulations 1950 in the UK. The Health and Safety Executive in the UK has produced guidelines on controlling exposure to respirable crystalline silica in potteries, and the British Ceramics Federation provide, as a free download, a guidance booklet. Environmental concerns include off-site water pollution, air pollution, disposal of hazardous materials, disposal of rejected ware and fuel consumption. History A great part of the history of pottery is prehistoric, part of past pre-literate cultures. Therefore, much of this history can only be found among the artifacts of archaeology. Because pottery is so durable, pottery and shards of pottery survive for millennia at archaeological sites, and are typically the most common and important type of artifact to survive. Many prehistoric cultures are named after the pottery that is the easiest way to identify their sites, and archaeologists develop the ability to recognise different types from the chemistry of small shards. Before pottery becomes part of a culture, several conditions must generally be met. First, there must be usable clay available. Archaeological sites where the earliest pottery was found were near deposits of readily available clay that could be properly shaped and fired. China has large deposits of a variety of clay, which gave them an advantage in early development of fine pottery. Many countries have large deposits of a variety of clay. Second, it must be possible to heat the pottery to temperatures that will achieve the transformation from raw clay to ceramic. Methods to reliably create fires hot enough to fire pottery did not develop until late in the development of cultures. Third, the potter must have time available to prepare, shape and fire the clay into pottery. Even after control of fire was achieved, humans did not seem to develop pottery until a sedentary life was achieved. It has been hypothesized that pottery was developed only after humans established agriculture, which led to permanent settlements. However, the oldest known pottery is from the Czech Republic and dates to 28,000 BC, at the height of the most recent ice age, long before the beginnings of agriculture. Fourth, there must be a sufficient need for pottery in order to justify the resources required for its production. ===Early pottery=== Methods of forming: Hand-shaping was the earliest method used to form vessels. This included the combination of pinching and coiling. Firing: The earliest method for firing pottery wares was the use of bonfires pit fired pottery. Firing times might be short but the peak-temperatures achieved in the fire could be high, perhaps in the region of , and were reached very quickly. Clay: Early potters used whatever clay was available to them in their geographic vicinity. However, the lowest quality common red clay was adequate for low-temperature fires used for the earliest pots. Clay tempered with sand, grit, crushed shell or crushed pottery were often used to make bonfire-fired ceramics because they provided an open-body texture that allowed water and volatile components of the clay to escape freely. The coarser particles in the clay also acted to restrain shrinkage during drying, and hence reduce the risk of cracking. Form: In the main, early bonfire-fired wares were made with rounded bottoms to avoid sharp angles that might be susceptible to cracking. Glazing: the earliest pots were not glazed. The potter's wheel was invented in Europe in the 5th millennium BC, and revolutionised pottery production. Earliest potter's wheel dated to the middle of the 5th millennium BC from the Cucuteni–Trypillia culture in western Ukraine. Moulds were used to a limited extent as early as the 5th and 6th century BC by the Etruscans and more extensively by the Romans. Slipcasting, a popular method for shaping irregular shaped articles. It was first practised, to a limited extent, in China as early as the Tang dynasty. Transition to kilns: The earliest intentionally constructed were pit-kilns or trench-kilns, holes dug in the ground and covered with fuel. Holes in the ground provided insulation and resulted in better control over firing. Kilns: Pit fire methods were adequate for simple earthenware, but other pottery types needed more sophisticated kilns. History by region Beginnings of pottery Pottery may well have been discovered independently in various places, probably by accidentally creating it at the bottom of fires on a clay soil. The earliest-known ceramic objects are Gravettian figurines such as those discovered at Dolní Věstonice in the modern-day Czech Republic. The Venus of Dolní Věstonice is a Venus figurine, a statuette of a nude female figure dated to 29,000–25,000 BC (Gravettian industry). But there is no evidence of pottery vessels from this period. Weights for looms or fishing-nets are a very common use for the earliest pottery. Sherds have been found in China and Japan from a period between 12,000 and perhaps as long as 18,000 years ago. As of 2012, the earliest pottery vessels found anywhere in the world, dating to 20,000 to 19,000 years before the present, was found at Xianren Cave in the Jiangxi province of China. Other early pottery vessels include those excavated from the Yuchanyan Cave in southern China, dated from 16,000 BC, and those found in the Amur River basin in the Russian Far East, dated from 14,000 BC. The Odai Yamamoto I site, belonging to the Jōmon period, currently has the oldest pottery in Japan. Excavations in 1998 uncovered earthenware fragments which have been dated as early as 14,500 BC. The term "Jōmon" means "cord-marked" in Japanese. This refers to the markings made on the vessels and figures using sticks with cords during their production. Recent research has elucidated how Jōmon pottery was used by its creators. It appears that pottery was independently developed in Sub-Saharan Africa during the 10th millennium BC, with findings dating to at least 9,400 BC from central Mali, and in South America during the 9,000s–7,000s BC. The Malian finds date to the same period as similar finds from East Asia – the triangle between Siberia, China and Japan – and are associated in both regions to the same climatic changes (at the end of the ice age new grassland develops, enabling hunter-gatherers to expand their habitat), met independently by both cultures with similar developments: the creation of pottery for the storage of wild cereals (pearl millet), and that of small arrowheads for hunting small game typical of grassland. Alternatively, the creation of pottery in the case of the Incipient Jōmon civilisation could be due to the intensive exploitation of freshwater and marine organisms by late glacial foragers, who started developing ceramic containers for their catch. East Asia In Japan, the Jōmon period has a long history of development of Jōmon pottery which was characterized by impressions of rope on the surface of the pottery created by pressing rope into the clay before firing. Glazed Stoneware was being created as early as the 15th century BC in China. A form of Chinese porcelain became a significant Chinese export from the Tang dynasty (AD 618–906) onwards. Korean potters adopted porcelain as early as the 14th century AD. The ceramic industry has developed greatly since the Goryeo dynasty, and Goryeo ware, a celadon with unique inlaying techniques, was produced. Later, when white porcelain became common and celadon fell, they created unique ceramics such as Buncheong. Japan's white porcelain was influenced by potters kidnapped during the Japanese invasions of Korea (1592–1598), called The Ceramic Wars, and Japanese engineers introduced it during the Fall of the Ming dynasty's. Typically, Korean potters who settled in Arita pass on pottery techniques, Shonzui Goradoyu-go brought back the secret of its manufacture from the Chinese kilns at Jingdezhen. In contrast to Europe, the Chinese social elite used pottery extensively at table, for religious purposes, and for decoration, and the standards of fine pottery were very high. From the Song dynasty (960–1279) for several centuries, the tastes of Chinese elites favoured plain-coloured and exquisitely formed pieces; during this period porcelain was perfected in Ding ware, although it was the only one of the Five Great Kilns of the Song period to use it. The traditional Chinese category of high-fired wares includes stoneware types such as Ru ware, Longquan celadon and Guan ware. Painted wares such as Cizhou ware had a lower status, though they were acceptable for making pillows. The arrival of Chinese blue and white porcelain was probably a product of the Mongol Yuan dynasty (1271–1368) dispersing artists and craftsmen across its large empire. Both the cobalt stains used for the blue colour, and the style of painted decoration, usually based on plant shapes, were initially borrowed from the Islamic world, which the Mongols had also conquered. At the same time Jingdezhen porcelain, produced in Imperial factories, took the undisputed leading role in production. The new elaborately painted style was now favoured at court, and gradually more colours were added. The secret of making such porcelain was sought in the Islamic world and later in Europe when examples were imported from the East. Many attempts were made to imitate it in Italy and France. However it was not produced outside of East Asia until 1709 in Germany. South Asia Cord-Impressed style pottery belongs to "Mesolithic" ceramic tradition that developed among Vindhya hunter-gatherers in Central India during the Mesolithic period. This ceramic style is also found in later Proto-Neolithic phase in nearby regions. This early type of pottery, also found at the site of Lahuradewa, is currently the oldest known pottery tradition in South Asia, dating back to 7,000–6,000 BC. Wheel-made pottery began to be made during the Mehrgarh Period II (5,500–4,800 BC) and Merhgarh Period III (4,800–3,500 BC), known as the ceramic Neolithic and chalcolithic. Pottery, including items known as the ed-Dur vessels, originated in regions of the Saraswati River / Indus River and have been found in a number of sites in the Indus Civilization. Despite an extensive prehistoric record of pottery, including painted wares, little "fine" or luxury pottery was made in the subcontinent in historic times. Hinduism discourages eating off pottery, which probably largely accounts for this. Most traditional Indian pottery vessels are large pots or jars for storage, or small cups or lamps, occasionally treated as disposable. In contrast there are long traditions of sculpted figures, often rather large, in terracotta; this continues with the Bankura horses in Panchmura, West Bengal. Southeast Asia Pottery in Southeast Asia is as diverse as its ethnic groups. Each ethnic group has their own set of standards when it comes to pottery arts. Potteries are made due to various reasons, such as trade, food and beverage storage, kitchen usage, religious ceremonies, and burial purposes. West Asia Around 8000 BC during the Pre-pottery Neolithic period, and before the invention of pottery, several early settlements became experts in crafting beautiful and highly sophisticated containers from stone, using materials such as alabaster or granite, and employing sand to shape and polish. Artisans used the veins in the material to maximum visual effect. Such objects have been found in abundance on the upper Euphrates river, in what is today eastern Syria, especially at the site of Bouqras. The earliest history of pottery production in the Fertile Crescent starts the Pottery Neolithic and can be divided into four periods, namely: the Hassuna period (7000–6500 BC), the Halaf period (6500–5500 BC), the Ubaid period (5500–4000 BC), and the Uruk period (4000–3100 BC). By about 5000 BC pottery-making was becoming widespread across the region, and spreading out from it to neighbouring areas. Pottery making began in the 7th millennium BC. The earliest forms, which were found at the Hassuna site, were hand formed from slabs, undecorated, unglazed low-fired pots made from reddish-brown clays. Within the next millennium, wares were decorated with elaborate painted designs and natural forms, incising and burnished. The invention of the potter's wheel in Mesopotamia sometime between 6,000 and 4,000 BC (Ubaid period) revolutionised pottery production. Newer kiln designs could fire wares to to which enabled increased possibilities. Production was now carried out by small groups of potters for small cities, rather than individuals making wares for a family. The shapes and range of uses for ceramics and pottery expanded beyond simple vessels to store and carry to specialized cooking utensils, pot stands and rat traps. As the region developed new organizations and political forms, pottery became more elaborate and varied. Some wares were made using moulds, allowing for increased production for the needs of the growing populations. Glazing was commonly used and pottery was more decorated. In the Chalcolithic period in Mesopotamia, Halafian pottery achieved a level of technical competence and sophistication, not seen until the later developments of Greek pottery with Corinthian and Attic ware. Europe Europe's oldest pottery, dating from circa 6700 BC, was found on the banks of the Samara River in the middle Volga region of Russia. These sites are known as the Yelshanka culture. The early inhabitants of Europe developed pottery in the Linear Pottery culture slightly later than the Near East, circa 5500–4500 BC. In the ancient Western Mediterranean elaborately painted earthenware reached very high levels of artistic achievement in the Greek world; there are large numbers of survivals from tombs. Minoan pottery was characterized by complex painted decoration with natural themes. The classical Greek culture began to emerge around 1000 BC featuring a variety of well crafted pottery which now included the human form as a decorating motif. The pottery wheel was now in regular use. Although glazing was known to these potters, it was not widely used. Instead, a more porous clay slip was used for decoration. A wide range of shapes for different uses developed early and remained essentially unchanged during Greek history. Fine Etruscan pottery was heavily influenced by Greek pottery and often imported Greek potters and painters. Ancient Roman pottery made much less use of painting, but used moulded decoration, allowing industrialized production on a huge scale. Much of the so-called red Samian ware of the Early Roman Empire was produced in modern Germany and France, where entrepreneurs established large potteries. Excavations at Augusta Raurica, near Basel, Switzerland, have revealed a pottery production site in use from the 1st to the 4th century AD. Pottery was hardly seen on the tables of elites from Hellenistic times until the Renaissance, and most medieval wares were coarse and utilitarian, as the elites ate off metal vessels. Painted Hispano-Moresque ware from Spain, developing the styles of Al-Andalus, became a luxury for late medieval elites, and was adapted in Italy into maiolica in the Italian Renaissance. Both of these were faience or tin-glazed earthenware, and fine faience continued to be made until around 1800 in various countries, especially France, with Nevers faience and several other centres. In the 17th century, imports of Chinese export porcelain and its Japanese equivalent raised the market expectations of fine pottery, and European manufacturers eventually learned to make porcelain, often in the form of soft-paste porcelain, and from the 18th century European porcelain and other wares from a great number of producers became extremely popular, reducing Asian imports. United Kingdom The city of Stoke-on-Trent is widely known as "The Potteries" because of the large number of pottery factories or, colloquially, "Pot Banks". It was one of the first industrial cities of the modern era where, as early as 1785, two hundred pottery manufacturers employed 20,000 workers. Josiah Wedgwood (1730–1795) was the dominant leader. In North Staffordshire hundreds of companies produced all kinds of pottery, from tablewares and decorative pieces to industrial items. The main pottery types of earthenware, stoneware and porcelain were all made in large quantities, and the Staffordshire industry was a major innovator in developing new varieties of ceramic bodies such as bone china and jasperware, as well as pioneering transfer printing and other glazing and decorating techniques. In general Staffordshire was strongest in the middle and low price ranges, though the finest and most expensive types of wares were also made. By the late 18th century North Staffordshire was the largest producer of ceramics in the UK, despite significant hubs elsewhere. Large export markets took Staffordshire pottery around the world, especially in the 19th century. Production had begun to decline in the late 19th century, as other countries developed their industries, and declined notably after World War II. Employment fell from 45,000 in 1975 to 23,000 in 1991, and 13,000 in 2002. Arabic pottery Early Islamic pottery followed the forms of the regions which the Arabs conquered. Eventually, however, there was cross-fertilization between the regions. This was most notable in the Chinese influences on Islamic pottery. Trade between China and Islam took place via the system of trading posts over the lengthy Silk Road. Middle Eastern nations imported stoneware and later porcelain from China. China imported the minerals for Cobalt blue from the Islamic ruled Persia to decorate their blue and white porcelain, which they then exported to the Islamic world. Likewise, Arabic art contributed to a lasting pottery form identified as Hispano-Moresque in Andalucia. Unique Islamic forms were also developed, including fritware, lusterware and specialized glazes like tin-glazing, which led to the development of the popular maiolica. One major emphasis in ceramic development in the Muslim world was the use of tile and decorative tilework. Americas Most evidence points to an independent development of pottery in the Native American cultures, with the earliest known dates from Brazil, from 9,500 to 5,000 years ago and 7,000 to 6,000 years ago. Further north in Mesoamerica, dates begin with the Archaic Era (3500–2000 BC), and into the Formative period (2000 BC – AD 200). These cultures did not develop the stoneware, porcelain or glazes found in the Old World. Maya ceramics include finely painted vessels, usually beakers, with elaborate scenes with several figures and texts. Several cultures, beginning with the Olmec, made terracotta sculpture, and sculptural pieces of humans or animals that are also vessels are produced in many places, with Moche portrait vessels among the finest. Africa The oldest pottery in the world outside of east Asia can be found in Africa. In 2007, Swiss archaeologists discovered pieces of some of the oldest pottery in Africa at Ounjougou in the central region of Mali, dating to at least 9,400 BC. Excavations in the Bosumpra Cave on the Kwahu Plateau in southeastern Ghana, have revealed well-manufactured pottery decorated with channelling and impressed peigne fileté rigide dating from the early tenth millennium cal. BC. Following the emergence of pottery traditions in the Ounjougou region of Mali around 11,900 BP and in the Bosumpra region of Ghana soon after, ceramics later arrived in the Iho Eleru region of Nigeria. In later periods, a relationship of the introduction of pot-making in some parts of Sub-Saharan Africa with the spread of Bantu languages has been long recognized, although the details remain controversial and awaiting further research, and no consensus has been reached. Oceania Pottery has been found in archaeological sites across the islands of Oceania. It is attributed to an ancient archaeological culture called the Lapita. Another form of pottery called Plainware is found throughout sites of Oceania. The relationship between Lapita pottery and Plainware is not altogether clear. The Indigenous Australians never developed pottery. After Europeans came to Australia and settled, they found deposits of clay which were analysed by English potters as excellent for making pottery. Less than 20 years later, Europeans came to Australia and began creating pottery. Since then, ceramic manufacturing, mass-produced pottery and studio pottery have flourished in Australia.	Technology	Visual arts	null
24651	https://en.wikipedia.org/wiki/Pantograph	Pantograph	A pantograph (, from their original use for copying writing) is a mechanical linkage connected in a manner based on parallelograms so that the movement of one pen, in tracing an image, produces identical movements in a second pen. If a line drawing is traced by the first point, an identical, enlarged, or miniaturized copy will be drawn by a pen fixed to the other. Using the same principle, different kinds of pantographs are used for other forms of duplication in areas such as sculpting, minting, engraving, and milling. History The ancient Greek engineer Hero of Alexandria described pantographs in his work Mechanics. In 1603, Christoph Scheiner used a pantograph to copy and scale diagrams, and wrote about the invention over 27 years later, in "Pantographice seu Ars delineandi res quaslibet per parallelogrammum lineare seu cavum" (Rome 1631). One arm of the pantograph contained a small pointer, while the other held a drawing implement, and by moving the pointer over a diagram, a copy of the diagram was drawn on another piece of paper. By changing the positions of the arms in the linkage between the pointer arm and drawing arm, the scale of the image produced can be changed. In 1821, Professor William Wallace (1768–1843) invented the eidograph to improve upon the practical utility of the pantograph. The eidograph relocates the fixed point to the center of the parallelogram and uses a narrow parallelogram to provide improved mechanical advantages. Uses Drafting The original use of the pantograph was for copying and scaling line drawings. Modern versions are sold as technical toys. Sculpture and minting Sculptors use a three-dimensional version of the pantograph, usually a large boom connected to a fixed point at one end, bearing two rotating pointing needles at arbitrary points along this boom. By adjusting the needles different enlargement or reduction ratios can be achieved. This device, now largely overtaken by computer guided router systems that scan a model and can produce it in a variety of materials and in any desired size, was invented by inventor and steam pioneer James Watt and perfected by Benjamin Cheverton in 1836. Cheverton's machine was fitted with a rotating cutting bit to carve reduced versions of well-known sculptures. A three-dimensional pantograph can also be used to enlarge sculpture by interchanging the position of the model and the copy. Another version is still very much in use to reduce the size of large relief designs for coins down to the required size of the coin. Acoustic cylinder duplication One advantage of phonograph and gramophone discs over cylinders in the 1890s—before electronic amplification was available—was that large numbers of discs could be stamped quickly and cheaply. In 1890, the only ways of manufacturing copies of a master cylinder were to mold the cylinders (which was slow and, early on, produced very poor copies), or to acoustically copy the sound by placing the horns of two phonographs together or to hook the two together with a rubber tube (one phonograph recording and the other playing the cylinder back). Instead of copying a master cylinder, the other alternative was to record a performance to multiple gramophones simultaneously, over and over again, making each cylinder a master copy. Edison, Bettini, Leon Douglass and others solved this problem (partly) by mechanically linking a cutting stylus and a playback stylus together and copying the "hill-and-dale" grooves of the cylinder mechanically. When molding improved somewhat, molded cylinders were used as pantograph masters. This was employed by Edison and Columbia in 1898, and was used until about January 1902 (Columbia brown waxes after this were molded). Some companies like the United States Phonograph Company of Newark, New Jersey, supplied cylinder masters for smaller companies so that they could duplicate them, sometimes pantographically. Pantographs could turn out about 30 records per day and produce up to about 150 records per master. In theory, pantograph masters could be used for 200 or 300 duplicates if the master and the duplicate were running in reverse and the record would be duplicated in reverse. This, in theory, could extend the usability of a pantograph master by using the unworn/lesser worn part of the recording for duplication. Pathé employed this system with mastering their vertically cut records until 1923; a , master cylinder, rotating at a high speed, would be recorded on. This was done as the resulting cylinder was considerably loud and of very high fidelity. Then, the cylinder would be placed on the mandrel of a duplicating pantograph that would be played with a stylus on the end of a lever, which would transfer the sound to a wax disc master, which would be electroplated and be used to stamp copies out. This system resulted in some fidelity reduction and rumble, but relatively high quality sound. Edison Diamond Disc Records were made by recording directly onto the wax master disc. Milling machines Before the advent of control technologies such as numerical control (NC and CNC) and programmable logic control (PLC), duplicate parts being milled on a milling machine could not have their contours mapped out by moving the milling cutter in a "connect-the-dots" ("by-the-numbers") fashion. The only ways to control the movement of the cutting tool were to dial the positions by hand using dexterous skill (with natural limits on a human's accuracy and precision) or to trace a cam, template, or model in some way, and have the cutter mimic the movement of the tracing stylus. If the milling head was mounted on a pantograph, a duplicate part could be cut (and at various scales of magnification besides 1:1) simply by tracing a template. (The template itself was usually made by a tool and die maker using toolroom methods, including milling via dialing followed by hand sculpting with files and/or die grinder points.) This was essentially the same concept as reproducing documents with a pen-equipped pantograph, but applied to the machining of hard materials such as metal, wood, or plastic. Pantograph routing, which is conceptually identical to pantograph milling, also exists (as does CNC routing). The Blanchard lathe, a copying lathe developed by Thomas Blanchard, used the same essential concept. The development and dissemination throughout industry of NC, CNC, PLC, and other control technologies provided a new way to control the movement of the milling cutter: via feeding information from a program to actuators (servos, selsyns, leadscrews, machine slides, spindles, and so on) that would move the cutter as the information directed. Today most commercial machining is done via such programmable, computerized methods. Home machinists are likely to work via manual control, but computerized control has reached the home-shop level as well (it is just not yet as pervasive as its commercial counterparts). Thus pantograph milling machines are largely a thing of the past. They are still in commercial use, but at a greatly reduced and ever-dwindling level. They are no longer built new by machine tool builders, but a small market for used machines still exists. As for the magnification-and-reduction feature of a pantograph (with the scale determined by the adjustable arm lengths), it is achieved in CNC via mathematic calculations that the computer applies to the program information practically instantaneously. Scaling functions (as well as mirroring functions) are built into languages such as G-code. Other uses In another application similar to drafting, the pantograph is incorporated into a pantograph engraving machine with a revolving cutter instead of a pen, and a tray at the pointer end to fix precut lettered plates (referred to as 'copy'), which the pointer follows and thus the cutter, via the pantograph, reproduces the 'copy' at a ratio to which the pantograph arms have been set. The typical range of ratio is Maximum 1:1 Minimum 50:1 (reduction) In this way machinists can neatly and accurately engrave numbers and letters onto a part. Pantographs are no longer commonly used in modern engraving, with computerized laser and rotary engraving taking favor. The device which maintains electrical contact with the contact wire and transfers power from the wire to the traction unit, used in electric locomotives and trams, is also called a "pantograph". Herman Hollerith's "Keyboard punch" used for the 1890 U.S. Census was a pantograph design and sometimes referred to as "The Pantograph Punch". An early 19th-century device employing this mechanism is the polygraph, which produces a duplicate of a letter as the original is written. In 1886, Eduard Selling patented a prize-winning calculating machine based on the pantograph, although it was not commercially successful. Longarm quilting machine operators may trace a pantograph, paper pattern, with a laser pointer to stitch a custom pattern onto the quilt. Digitized pantographs are followed by computerized machines. Linn Boyd Benton invented a pantographic engraving machine for type design, which was capable not only of scaling a single font design pattern to a variety of sizes, but could also condense, extend, and slant the design (mathematically, these are cases of affine transformation, which is the fundamental geometric operation of most systems of digital typography today, including PostScript). Richard Feynman used the analogy of a pantograph as a way of scaling down tools to the nanometer scale in his talk "There's Plenty of Room at the Bottom".	Technology	Mechanisms	null
24669	https://en.wikipedia.org/wiki/Pauli%20exclusion%20principle	Pauli exclusion principle	In quantum mechanics, the Pauli exclusion principle (German: Pauli-Ausschlussprinzip) states that two or more identical particles with half-integer spins (i.e. fermions) cannot simultaneously occupy the same quantum state within a system that obeys the laws of quantum mechanics. This principle was formulated by Austrian physicist Wolfgang Pauli in 1925 for electrons, and later extended to all fermions with his spin–statistics theorem of 1940. In the case of electrons in atoms, the exclusion principle can be stated as follows: in a poly-electron atom it is impossible for any two electrons to have the same two values of all four of their quantum numbers, which are: n, the principal quantum number; , the azimuthal quantum number; m, the magnetic quantum number; and ms, the spin quantum number. For example, if two electrons reside in the same orbital, then their values of n, , and m are equal. In that case, the two values of ms (spin) pair must be different. Since the only two possible values for the spin projection ms are +1/2 and −1/2, it follows that one electron must have ms = +1/2 and one ms = −1/2. Particles with an integer spin (bosons) are not subject to the Pauli exclusion principle. Any number of identical bosons can occupy the same quantum state, such as photons produced by a laser, or atoms found in a Bose–Einstein condensate. A more rigorous statement is: under the exchange of two identical particles, the total (many-particle) wave function is antisymmetric for fermions and symmetric for bosons. This means that if the space and spin coordinates of two identical particles are interchanged, then the total wave function changes sign for fermions, but does not change sign for bosons. So, if hypothetically two fermions were in the same statefor example, in the same atom in the same orbital with the same spinthen interchanging them would change nothing and the total wave function would be unchanged. However, the only way a total wave function can both change sign (required for fermions), and also remain unchanged is that such a function must be zero everywhere, which means such a state cannot exist. This reasoning does not apply to bosons because the sign does not change. Overview The Pauli exclusion principle describes the behavior of all fermions (particles with half-integer spin), while bosons (particles with integer spin) are subject to other principles. Fermions include elementary particles such as quarks, electrons and neutrinos. Additionally, baryons such as protons and neutrons (subatomic particles composed from three quarks) and some atoms (such as helium-3) are fermions, and are therefore described by the Pauli exclusion principle as well. Atoms can have different overall spin, which determines whether they are fermions or bosons: for example, helium-3 has spin 1/2 and is therefore a fermion, whereas helium-4 has spin 0 and is a boson. The Pauli exclusion principle underpins many properties of everyday matter, from its large-scale stability to the chemical behavior of atoms. Half-integer spin means that the intrinsic angular momentum value of fermions is (reduced Planck constant) times a half-integer (1/2, 3/2, 5/2, etc.). In the theory of quantum mechanics, fermions are described by antisymmetric states. In contrast, particles with integer spin (bosons) have symmetric wave functions and may share the same quantum states. Bosons include the photon, the Cooper pairs which are responsible for superconductivity, and the W and Z bosons. Fermions take their name from the Fermi–Dirac statistical distribution, which they obey, and bosons take theirs from the Bose–Einstein distribution. History In the early 20th century it became evident that atoms and molecules with even numbers of electrons are more chemically stable than those with odd numbers of electrons. In the 1916 article "The Atom and the Molecule" by Gilbert N. Lewis, for example, the third of his six postulates of chemical behavior states that the atom tends to hold an even number of electrons in any given shell, and especially to hold eight electrons, which he assumed to be typically arranged symmetrically at the eight corners of a cube. In 1919 chemist Irving Langmuir suggested that the periodic table could be explained if the electrons in an atom were connected or clustered in some manner. Groups of electrons were thought to occupy a set of electron shells around the nucleus. In 1922, Niels Bohr updated his model of the atom by assuming that certain numbers of electrons (for example 2, 8 and 18) corresponded to stable "closed shells". Pauli looked for an explanation for these numbers, which were at first only empirical. At the same time he was trying to explain experimental results of the Zeeman effect in atomic spectroscopy and in ferromagnetism. He found an essential clue in a 1924 paper by Edmund C. Stoner, which pointed out that, for a given value of the principal quantum number (n), the number of energy levels of a single electron in the alkali metal spectra in an external magnetic field, where all degenerate energy levels are separated, is equal to the number of electrons in the closed shell of the noble gases for the same value of n. This led Pauli to realize that the complicated numbers of electrons in closed shells can be reduced to the simple rule of one electron per state if the electron states are defined using four quantum numbers. For this purpose he introduced a new two-valued quantum number, identified by Samuel Goudsmit and George Uhlenbeck as electron spin. Connection to quantum state symmetry In his Nobel lecture, Pauli clarified the importance of quantum state symmetry to the exclusion principle: Among the different classes of symmetry, the most important ones (which moreover for two particles are the only ones) are the symmetrical class, in which the wave function does not change its value when the space and spin coordinates of two particles are permuted, and the antisymmetrical class, in which for such a permutation the wave function changes its sign...[The antisymmetrical class is] the correct and general wave mechanical formulation of the exclusion principle. The Pauli exclusion principle with a single-valued many-particle wavefunction is equivalent to requiring the wavefunction to be antisymmetric with respect to exchange. If and range over the basis vectors of the Hilbert space describing a one-particle system, then the tensor product produces the basis vectors of the Hilbert space describing a system of two such particles. Any two-particle state can be represented as a superposition (i.e. sum) of these basis vectors: where each is a (complex) scalar coefficient. Antisymmetry under exchange means that . This implies when , which is Pauli exclusion. It is true in any basis since local changes of basis keep antisymmetric matrices antisymmetric. Conversely, if the diagonal quantities are zero in every basis, then the wavefunction component is necessarily antisymmetric. To prove it, consider the matrix element This is zero, because the two particles have zero probability to both be in the superposition state . But this is equal to The first and last terms are diagonal elements and are zero, and the whole sum is equal to zero. So the wavefunction matrix elements obey: or For a system with particles, the multi-particle basis states become n-fold tensor products of one-particle basis states, and the coefficients of the wavefunction are identified by n one-particle states. The condition of antisymmetry states that the coefficients must flip sign whenever any two states are exchanged: for any . The exclusion principle is the consequence that, if for any then This shows that none of the n particles may be in the same state. Advanced quantum theory According to the spin–statistics theorem, particles with integer spin occupy symmetric quantum states, and particles with half-integer spin occupy antisymmetric states; furthermore, only integer or half-integer values of spin are allowed by the principles of quantum mechanics. In relativistic quantum field theory, the Pauli principle follows from applying a rotation operator in imaginary time to particles of half-integer spin. In one dimension, bosons, as well as fermions, can obey the exclusion principle. A one-dimensional Bose gas with delta-function repulsive interactions of infinite strength is equivalent to a gas of free fermions. The reason for this is that, in one dimension, the exchange of particles requires that they pass through each other; for infinitely strong repulsion this cannot happen. This model is described by a quantum nonlinear Schrödinger equation. In momentum space, the exclusion principle is valid also for finite repulsion in a Bose gas with delta-function interactions, as well as for interacting spins and Hubbard model in one dimension, and for other models solvable by Bethe ansatz. The ground state in models solvable by Bethe ansatz is a Fermi sphere. Applications Atoms The Pauli exclusion principle helps explain a wide variety of physical phenomena. One particularly important consequence of the principle is the elaborate electron shell structure of atoms and the way atoms share electrons, explaining the variety of chemical elements and their chemical combinations. An electrically neutral atom contains bound electrons equal in number to the protons in the nucleus. Electrons, being fermions, cannot occupy the same quantum state as other electrons, so electrons have to "stack" within an atom, i.e. have different spins while at the same electron orbital as described below. An example is the neutral helium atom (He), which has two bound electrons, both of which can occupy the lowest-energy (1s) states by acquiring opposite spin; as spin is part of the quantum state of the electron, the two electrons are in different quantum states and do not violate the Pauli principle. However, the spin can take only two different values (eigenvalues). In a lithium atom (Li), with three bound electrons, the third electron cannot reside in a 1s state and must occupy a higher-energy state instead. The lowest available state is 2s, so that the ground state of Li is 1s22s. Similarly, successively larger elements must have shells of successively higher energy. The chemical properties of an element largely depend on the number of electrons in the outermost shell; atoms with different numbers of occupied electron shells but the same number of electrons in the outermost shell have similar properties, which gives rise to the periodic table of the elements. To test the Pauli exclusion principle for the helium atom, Gordon Drake carried out very precise calculations for hypothetical states of the He atom that violate it, which are called paronic states. Later, K. Deilamian et al. used an atomic beam spectrometer to search for the paronic state 1s2s 1S0 calculated by Drake. The search was unsuccessful and showed that the statistical weight of this paronic state has an upper limit of . (The exclusion principle implies a weight of zero.) Solid state properties In conductors and semiconductors, there are very large numbers of molecular orbitals which effectively form a continuous band structure of energy levels. In strong conductors (metals) electrons are so degenerate that they cannot even contribute much to the thermal capacity of a metal. Many mechanical, electrical, magnetic, optical and chemical properties of solids are the direct consequence of Pauli exclusion. Stability of matter The stability of each electron state in an atom is described by the quantum theory of the atom, which shows that close approach of an electron to the nucleus necessarily increases the electron's kinetic energy, an application of the uncertainty principle of Heisenberg. However, stability of large systems with many electrons and many nucleons is a different question, and requires the Pauli exclusion principle. It has been shown that the Pauli exclusion principle is responsible for the fact that ordinary bulk matter is stable and occupies volume. This suggestion was first made in 1931 by Paul Ehrenfest, who pointed out that the electrons of each atom cannot all fall into the lowest-energy orbital and must occupy successively larger shells. Atoms, therefore, occupy a volume and cannot be squeezed too closely together. The first rigorous proof was provided in 1967 by Freeman Dyson and Andrew Lenard (de), who considered the balance of attractive (electron–nuclear) and repulsive (electron–electron and nuclear–nuclear) forces and showed that ordinary matter would collapse and occupy a much smaller volume without the Pauli principle. A much simpler proof was found later by Elliott H. Lieb and Walter Thirring in 1975. They provided a lower bound on the quantum energy in terms of the Thomas-Fermi model, which is stable due to a theorem of Teller. The proof used a lower bound on the kinetic energy which is now called the Lieb–Thirring inequality. The consequence of the Pauli principle here is that electrons of the same spin are kept apart by a repulsive exchange interaction, which is a short-range effect, acting simultaneously with the long-range electrostatic or Coulombic force. This effect is partly responsible for the everyday observation in the macroscopic world that two solid objects cannot be in the same place at the same time. Astrophysics Dyson and Lenard did not consider the extreme magnetic or gravitational forces that occur in some astronomical objects. In 1995 Elliott Lieb and coworkers showed that the Pauli principle still leads to stability in intense magnetic fields such as in neutron stars, although at a much higher density than in ordinary matter. It is a consequence of general relativity that, in sufficiently intense gravitational fields, matter collapses to form a black hole. Astronomy provides a spectacular demonstration of the effect of the Pauli principle, in the form of white dwarf and neutron stars. In both bodies, the atomic structure is disrupted by extreme pressure, but the stars are held in hydrostatic equilibrium by degeneracy pressure, also known as Fermi pressure. This exotic form of matter is known as degenerate matter. The immense gravitational force of a star's mass is normally held in equilibrium by thermal pressure caused by heat produced in thermonuclear fusion in the star's core. In white dwarfs, which do not undergo nuclear fusion, an opposing force to gravity is provided by electron degeneracy pressure. In neutron stars, subject to even stronger gravitational forces, electrons have merged with protons to form neutrons. Neutrons are capable of producing an even higher degeneracy pressure, neutron degeneracy pressure, albeit over a shorter range. This can stabilize neutron stars from further collapse, but at a smaller size and higher density than a white dwarf. Neutron stars are the most "rigid" objects known; their Young modulus (or more accurately, bulk modulus) is 20 orders of magnitude larger than that of diamond. However, even this enormous rigidity can be overcome by the gravitational field of a neutron star mass exceeding the Tolman–Oppenheimer–Volkoff limit, leading to the formation of a black hole.	Physical sciences	Quantum mechanics	null
24710	https://en.wikipedia.org/wiki/North%20American%20P-51%20Mustang	North American P-51 Mustang	The North American Aviation P-51 Mustang is an American long-range, single-seat fighter and fighter-bomber used during World War II and the Korean War, among other conflicts. The Mustang was designed in 1940 by a team headed by James H. Kindelberger of North American Aviation (NAA) in response to a requirement of the British Purchasing Commission. The commission approached NAA to build Curtiss P-40 fighters under license for the Royal Air Force (RAF). Rather than build an old design from another company, NAA proposed the design and production of a more modern fighter. The prototype NA-73X airframe was completed on 9 September 1940, 102 days after contract signing, achieving its first flight on 26 October. The Mustang was designed to use the Allison V-1710 engine without an export-sensitive turbosupercharger or a multi-stage supercharger, resulting in limited high-altitude performance. The aircraft was first flown operationally and very successfully by the RAF and as a tactical-reconnaissance aircraft and fighter-bomber (Mustang Mk I). In mid 1942, a development project known as the Rolls-Royce Mustang X, replaced the Allison engine with a Rolls-Royce Merlin 65 two-stage inter-cooled supercharged engine. During testing at Rolls-Royce's airfield at Hucknall in England, it was clear the engine dramatically improved the aircraft's performance at altitudes above without sacrificing range. Following receipt of the test results and after further flights by a number of USAAF pilots, the results were so positive that North American began work on converting several aircraft developing into the P-51B/C (Mustang Mk III) model, which became the first long range fighter to be able to compete with the Luftwaffe's fighters. The definitive version, the P-51D, was powered by the Packard V-1650-7, a license-built version of the two-speed, two-stage-supercharged Merlin 66, and was armed with six .50 caliber (12.7 mm) AN/M2 Browning machine guns. From late 1943, P-51Bs and P-51Cs (supplemented by P-51Ds from mid-1944) were used by the USAAF's Eighth Air Force to escort bombers in raids over Germany, while the RAF's Second Tactical Air Force and the USAAF's Ninth Air Force used the Merlin-powered Mustangs as fighter-bombers, roles in which the Mustang helped ensure Allied air superiority in 1944. The P-51 was also used by Allied air forces in the North African, Mediterranean, Italian, and Pacific theaters. During World War II, Mustang pilots claimed to have destroyed 4,950 enemy aircraft. At the start of the Korean War, the Mustang, by then redesignated F-51, was the main fighter of the United States until jet fighters, including North American's F-86 Sabre, took over this role; the Mustang then became a specialized fighter-bomber. Despite the advent of jet fighters, the Mustang remained in service with some air forces until the early 1980s. After the Korean War, Mustangs became popular civilian warbirds and air racing aircraft. Design and development In 1938, the British government established a purchasing commission in the United States, headed by Sir Henry Self. Self was given overall responsibility for RAF production, research, and development, and also served with Sir Wilfrid Freeman, the Air Member for Development and Production. Self also sat on the British Air Council Subcommittee on Supply (or "Supply Committee"), and one of his tasks was to organize the manufacturing and supply of American fighter aircraft for the RAF. At the time, the choice was very limited, as no U.S. aircraft then in production or flying met European standards, with only the Curtiss P-40 Tomahawk coming close. The Curtiss-Wright plant was running at capacity, so P-40s were in short supply. North American Aviation (NAA) was already supplying its T-6 Texan (known in British service as the "Harvard") trainer to the RAF, but was otherwise underused. NAA President "Dutch" Kindelberger approached Self to sell a new medium bomber, the North American B-25 Mitchell. Instead, Self asked if NAA could manufacture P-40s under license from Curtiss. Kindelberger said NAA could have a better aircraft with the same Allison V-1710 engine in the air sooner than establishing a production line for the P-40. John Attwood of NAA spent much time from January to April 1940 at the British Purchasing Commission's offices in New York discussing the British specifications of the proposed aircraft with British engineers. The discussions consisted of free-hand conceptual drawings of an aircraft with the British officials. Self was concerned that NAA had not ever designed a fighter, insisting they obtain the drawings and study the wind-tunnel test results for the P-40, before presenting them with detailed design drawings based on the agreed concept. NAA purchased the drawings and data from Curtiss for £56,000, confirming the purchase with the British Purchasing Commission. The commission approved the resulting detailed design drawings, signing the commencement of the Mustang project on 4 May 1940, and firmly ordering 320 on 29 May 1940. Prior to this, NAA only had a letter of intent for an order of 320 aircraft. Curtiss engineers accused NAA of plagiarism. The British Purchasing Commission stipulated armament of four .303 in (7.7 mm) machine guns (as used on the Tomahawk), a unit cost of no more than $40,000, and delivery of the first production aircraft by January 1941. In March 1940, 320 aircraft were ordered by Freeman, who had become the executive head of the Ministry of Aircraft Production (MAP) and the contract was promulgated on 24 April. The NA-73X, which was designed by a team led by lead engineer Edgar Schmued, followed the best conventional practice of the era, designed for ease of mass manufacturing. The design included several new features. One was a wing designed using laminar flow airfoils, which were developed co-operatively by NAA and the National Advisory Committee for Aeronautics (NACA). These airfoils generated low drag at high speeds. During the development of the NA-73X, a wind-tunnel test of two wings, one using NACA five-digit airfoils and the other using the new NAA/NACA 45–100 airfoils, was performed in the University of Washington Kirsten Wind Tunnel. The results of this test showed the superiority of the wing designed with the NAA/NACA 45–100 airfoils. The other feature was a new cooling arrangement positioned aft (single ducted water and oil radiators assembly) that reduced the fuselage drag and effects on the wing. Later, after much development, they discovered that the cooling assembly could take advantage of the Meredith effect, in which heated air exited the radiator with a slight amount of jet thrust. Because NAA lacked a suitable wind tunnel to test this feature, it used the GALCIT wind tunnel at the California Institute of Technology. This led to some controversy over whether the Mustang's cooling system aerodynamics were developed by NAA's engineer Schmued or by Curtiss, as NAA had purchased the complete set of P-40 wind tunnel data and flight test reports. The NA-73X was also one of the first aircraft to have a fuselage lofted mathematically using conic sections; this resulted in smooth, low-drag surfaces. To aid production, the airframe was divided into five main sections—forward, center, rear fuselage, and two wing halves—all of which were fitted with wiring and piping before being joined. The prototype NA-73X was rolled out in September 1940, just 102 days after the order had been placed; it first flew on 26 October 1940, 149 days into the contract, an uncommonly short development period even during the war. With test pilot Vance Breese at the controls, the prototype handled well and accommodated an impressive fuel load. The aircraft's three-section, semi-monocoque fuselage was constructed entirely of 24S aluminum alloy (a type of Duralumin) to save weight. It was armed with four .30 caliber (7.62 mm) AN/M2 Browning machine guns in the wings and two .50 caliber (12.7 mm) AN/M2 Browning machine guns mounted under the engine and firing through the propeller arc using a gun-synchronizing gear. While the USAAC could block any sales it considered detrimental to the interests of the US, the NA-73 was considered to be a special case because it had been designed at the behest of the British and all dealings were directly between the BPC and NAA, and did not involve the US Army or Wright Field in any way. In September 1940, a further 300 NA-73s were ordered by the MAP. To ensure uninterrupted delivery, Colonel Oliver P. Echols arranged with the Anglo-French Purchasing Commission to deliver the aircraft and NAA gave two examples (41-038 and 41-039) to the USAAC for evaluation. It is important to note that the Mustang I (NA-73 and NA-83) and the Ia (NA-91), produced for the British, were not equivalent to the P-51A which was a later model (NA-99). Two British Mustang Is were held back by the USAAF and given the provisional model number XP-51. The USAAF held back 57 Mustang Ia aircraft armed with 4 x 20mm Hispano cannon, from the third British order, converting most of them to tactical reconnaissance aircraft and designating them P-51-2/F6A. North American retained the second aircraft of this batch to help develop the P-51A. The Allison engine in the Mustang I had a single-stage supercharger that caused power to drop off rapidly above . This made it unsuitable for use at the altitudes where combat was taking place in Europe. Allison's attempts at developing a high-altitude engine were underfunded, but produced the V-1710-45, which featured a variable-speed auxiliary supercharger and developed at . In November 1941, NAA studied the possibility of using it, but fitting its excessive length in the Mustang would require extensive airframe modifications and cause long production delays. In May 1942, following positive reports from the RAF on the Mustang I's performance below 15,000 ft, Ronald Harker, a test pilot for Rolls-Royce, suggested fitting a Merlin 61, as fitted to the Spitfire Mk IX. The Merlin 61 had a two-speed, two-stage, intercooled supercharger, designed by Stanley Hooker of Rolls-Royce. Both the Merlin 61 and V-1710-39 were capable of about war emergency power at relatively low altitudes, but the Merlin developed at versus the Allison's at , delivering an increase in top speed from at ~ to an estimated at . In the end the Merlin 61 was never fitted to the Mustang X, (or any other Mustang). The 65 series (a medium altitude engine) was fitted to all Mustang X prototypes. Initially, the Mustang's steadfast champion, USAAC/F Assistant Air Attaché Major Thomas Hitchcock, was concerned that the USAAF had little or no interest in the potential of the P-51A and its development with the Merlin engine. He wrote: "Its development in this theatre has suffered for various reasons. Sired by the English out of an American mother, the Mustang has no parent in the Army Air Corps to appreciate and push its good points. It does not fully satisfy good people on both sides of the Atlantic who seem more interested in pointing with pride to the development of a 100% national product..." Nevertheless, during the British service development program of the Mustang I at Rolls-Royce's airfield at Hucknall, a close relationship was developed between NAA, the RAF Air Fighting Development Unit and Rolls Royce Rolls-Royce Flight Test Establishment at Hucknall. Following extensive communication between Hitchcock (based in England), Rolls Royce engineers and Phillip Legarra at NAA regarding the promising outlook of a Merlin Mustang, along with the subsequent work in progress by Rolls Royce on the Mustang X, NAA representatives including Mustang designer Schmued visited the UK to examine and discuss the project in detail. The promising calculations and modification progress by Rolls Royce led in July 1942 to a contract being let for two NAA Merlin prototypes, briefly designated XP-78, but soon to become the XP-51B. Based on the Packard V-1650-3 duplicating the Merlin 61's performance, NAA estimated for the XP-78 a top speed of at , and a service ceiling of . Initial flights of what was known to Rolls-Royce as the Mustang X were completed at Hucknall in October 1942. The first flight of the US version, designated XP-51B took place in November 1942, but the USAAF had become so interested in the Merlin Mustang project that an initial contract for 400 aircraft was placed three months beforehand in August. The conversion led to production of the P-51B beginning at NAA's Inglewood, California, plant in June 1943, and P-51s started to become available to the 8th and 9th air forces in the winter of 1943–1944. Conversion to the two-stage supercharged and intercooled Merlin 60 series, over heavier than the single-stage Allison, driving a four-bladed Hamilton Standard propeller, required moving the wing slightly forward to correct the aircraft's center of gravity. After the USAAF, in July 1943, directed fighter aircraft manufacturers to maximize internal fuel capacity, NAA calculated the P-51B's center of gravity to be forward enough to include an additional fuel tank in the fuselage behind the pilot, greatly increasing the aircraft's range over that of the earlier P-51A. NAA incorporated the tank in the production of the P-51B-10, and supplied kits to retrofit it to all existing P-51Bs. Operational history United Kingdom operational service The Mustang was initially developed for the RAF, which was its first user. As the first Mustangs were built to British requirements, these aircraft used factory numbers and were not P-51s; the order comprised 320 NA-73s, followed by 300 NA-83s, all of which were designated Mustang Mark I by the RAF. The first RAF Mustangs supplied under Lend-Lease were 93 Mk Ia designated as P-51s by the USAAF, followed by 50 P-51As used as Mustang Mk IIs. Aircraft supplied to Britain under Lend-Lease were required for accounting purposes to be on the USAAC's books before they could be supplied to Britain, but the British Aircraft Purchasing Commission signed its first contract for the North American NA-73 on 24 April 1940, before Lend-Lease was in effect. Thus, the initial order for the P-51 Mustang (as it was later known) was placed by the British under the "cash and carry" program, as required by the US Neutrality Acts of the 1930s. After the arrival of the initial aircraft in the UK in October 1941, the first squadron of Mustang Mk Is entered service in January 1942, the first being No. 26 Squadron RAF. Due to poor high-altitude performance, the Mustangs were used by Army Co-operation Command, rather than Fighter Command, and were used for tactical reconnaissance and ground-attack duties. On 10 May 1942, Mustangs first flew over France, near Berck-sur-Mer. On 27 July 1942, 16 RAF Mustangs undertook their first long-range reconnaissance mission over Germany. During the amphibious Dieppe Raid on the French coast (19 August 1942), four British and Canadian Mustang squadrons, including 26 Squadron, saw action covering the assault on the ground. By 1943–1944, British Mustangs were used extensively to seek out V-1 flying bomb sites. The last RAF Mustang Mk I and Mustang Mk II aircraft were struck off charge in 1945. Army Co-operation Command used the Mustang's superior speed and long range to conduct low-altitude "Rhubarb" raids over continental Europe, sometimes penetrating German airspace. The V-1710 engine ran smoothly at 1,100 rpm, versus 1,600 for the Merlin, enabling long flights over water at altitude before approaching the enemy coastline. Over land, these flights followed a zig-zag course, turning every six minutes to foil enemy attempts at plotting an interception. During the first 18 months of Rhubarb raids, RAF Mustang Mk.Is and Mk.Ias destroyed or heavily damaged 200 locomotives, over 200 canal barges, and an unknown number of enemy aircraft parked on the ground, for a loss of eight Mustangs. At sea level, the Mustangs were able to outrun all enemy aircraft encountered. The RAF gained a significant performance enhancement at low altitude by removing or resetting the engine's manifold pressure regulator to allow overboosting, raising output as high as 1,780 horsepower at 70 in Hg. In December 1942, Allison approved only 1,570 horsepower at 60 in Hg manifold pressure for the V-1710-39. The RAF later operated 308 P-51Bs and 636 P-51Cs, which were known in RAF service as Mustang Mk IIIs; the first units converted to the type in late 1943 and early 1944. Mustang Mk III units were operational until the end of World War II, though many units had already converted to the Mustang Mk IV (P-51D) and Mk IVa (P-51K) (828 in total, comprising 282 Mk IV and 600 Mk IVa). As all except the earliest aircraft were obtained under Lend-Lease, all Mustang aircraft still on RAF charge at the end of the war were either returned to the USAAF "on paper" or retained by the RAF for scrapping. The last RAF Mustangs were retired from service in 1947. U.S. operational service Prewar theory Prewar doctrine was based on the idea "the bomber will always get through". Despite RAF and Luftwaffe experience with daylight bombing, the USAAF still incorrectly believed in 1942 that tightly packed formations of bombers would have so much firepower that they could fend off fighters on their own. Fighter escort was a low priority, but when the concept was discussed in 1941, the Lockheed P-38 Lightning was considered to be most appropriate, as it had the speed and range. Another school of thought favored a heavily up-armed "gunship" conversion of a strategic bomber. A single-engined, high-speed fighter with the range of a bomber was thought to be an engineering impossibility. Eighth Air Force bomber operations 1942–1943 The 8th Air Force started operations from Britain in August 1942. At first, because of the limited scale of operations, no conclusive evidence showed American doctrine was failing. In the 26 operations flown to the end of 1942, the loss rate had been under 2%. In January 1943, at the Casablanca Conference, the Allies formulated the Combined Bomber Offensive (CBO) plan for "round-the-clock" bombing – USAAF daytime operations complementing the RAF nighttime raids on industrial centers. In June 1943, the Combined Chiefs of Staff issued the Pointblank Directive to destroy the Luftwaffe's capacity before the planned invasion of Europe, putting the CBO into full implementation. German daytime fighter efforts were, at that time, focused on the Eastern Front and several other distant locations. Initial efforts by the 8th met limited and unorganized resistance, but with every mission, the Luftwaffe moved more aircraft to the west and quickly improved their battle direction. In fall 1943, the 8th Air Force's heavy bombers conducted a series of deep-penetration raids into Germany, beyond the range of escort fighters. The Schweinfurt–Regensburg mission in August lost 60 B-17s of a force of 376, the 14 October attack lost 77 of a force of 291—26% of the attacking force. For the US, the very concept of self-defending bombers was called into question, but instead of abandoning daylight raids and turning to night bombing, as the RAF suggested, they chose other paths; at first, bombers converted to gunships (the Boeing YB-40) were believed to be able to escort the bomber formations, but when the concept proved to be unsuccessful, thoughts then turned to the Lockheed P-38 Lightning. In early 1943, the USAAF also decided that the Republic P-47 Thunderbolt and P-51B be considered for the roles of smaller escort fighters, and in July, a report stated that the P-51B was "the most promising plane" with an endurance of 4 hours 45 minutes with the standard internal fuel of 184 gallons plus 150 gallons carried externally. In August, a P-51B was fitted with an extra internal 85-gallon tank, but problems with longitudinal stability occurred, so some compromises in performance with the full tank were made. Since the fuel from the fuselage tank was used during the initial stages of a mission, the fuel tank would be fitted in all Mustangs destined for VIII Fighter Command. P-51 introduction The P-51 Mustang was a solution to the need for an effective bomber escort. It used a common, reliable engine and had internal space for a larger-than-average fuel load. With external fuel tanks, it could accompany the bombers from England to Germany and back. By the time the Pointblank offensive resumed in early 1944, matters had changed. Bomber escort defenses were initially layered, using the shorter-range P-38s and P-47s to escort the bombers during the initial stages of the raid before handing over to the P-51s when they were forced to turn for home. This provided continuous coverage during the raid. The Mustang was so clearly superior to earlier US designs that the 8th Air Force began to steadily switch its fighter groups to the Mustang, first swapping arriving P-47 groups to the 9th Air Force in exchange for those that were using P-51s, then gradually converting its Thunderbolt and Lightning groups. By the end of 1944, 14 of its 15 groups flew Mustangs. The Luftwaffe's twin-engined Messerschmitt Bf 110 heavy fighters brought up to deal with the bombers proved to be easy prey for the Mustangs, and had to be quickly withdrawn from combat. The Focke-Wulf Fw 190A, already suffering from poor high-altitude performance, was outperformed by the Mustang at the B-17's altitude, and when laden with heavy bomber-hunting weapons as a replacement for the more vulnerable twin-engined Zerstörer heavy fighters, it suffered heavy losses. The Messerschmitt Bf 109 had comparable performance at high altitudes, but its lightweight airframe was even more greatly affected by increases in armament. The Mustang's much lighter armament, tuned for antifighter combat, allowed it to overcome these single-engined opponents. Fighting the Luftwaffe At the start of 1944, Major General James Doolittle, the new commander of the 8th Air Force, released most fighters from the requirement of flying in close formation with the bombers, allowing them free rein to attack the Luftwaffe wherever it could be found. The aim was to achieve air supremacy. Mustang groups were sent far ahead of the bombers in a "fighter sweep" to intercept German fighters. Bomber crews complained, but by June, supremacy was achieved. The Luftwaffe answered with the Gefechtsverband ("battle formation"). This consisted of a Sturmgruppe of heavily armed and armored Fw 190As escorted by two Begleitgruppen of Bf 109s, whose task was to keep the Mustangs away from the Fw 190s as they attacked the bombers. This strategy proved to be problematic, as the large German formation took a long time to assemble and was difficult to maneuver. It was often intercepted by the P-51 "fighter sweeps" before it could attack the bombers. However, German attacks against bombers could be effective when they did occur; the bomber-destroyer Fw 190As swept in from astern and often pressed their attacks to within . While not always able to avoid contact with the escorts, the threat of mass attacks and later the "company front" (eight abreast) assaults by armored Sturmgruppe Fw 190As brought an urgency to attacking the Luftwaffe wherever it could be found, either in the air or on the ground. Beginning in late February 1944, 8th Air Force fighter units began systematic strafing attacks on German airfields with increasing frequency and intensity throughout the spring, with the objective of gaining air supremacy over the Normandy battlefield. In general, these were conducted by units returning from escort missions, but beginning in March, many groups also were assigned airfield attacks instead of bomber support. The P-51, particularly with the advent of the K-14 gyro gunsight and the development of "Clobber Colleges" for the training of fighter pilots in fall 1944, was a decisive element in Allied countermeasures against the Jagdverbände. The numerical superiority of the USAAF fighters, superb flying characteristics of the P-51, and pilot proficiency helped cripple the Luftwaffes fighter force. As a result, the fighter threat to the US, and later British, bombers was greatly diminished by July 1944. The RAF, long proponents of night bombing for protection, were able to reopen daylight bombing in 1944 as a result of the crippling of the Luftwaffe fighter arm. Reichsmarschall Hermann Göring, commander of the Luftwaffe during the war, was quoted as saying, "When I saw Mustangs over Berlin, I knew the jig was up." Beyond Pointblank On 15 April 1944, VIII Fighter Command began "Operation Jackpot", attacks on Luftwaffe fighter airfields. As the efficacy of these missions increased, the number of fighters at the German airbases fell to the point where they were no longer considered worthwhile targets. On 21 May, targets were expanded to include railways, locomotives, and other rolling stock used by the Germans to transport materiel and troops, in missions dubbed "Chattanooga". The P-51 excelled at this mission, although losses were much higher on strafing missions than in air-to-air combat, partially because the Mustang's liquid-cooled engine (particularly its liquid coolant system) was vulnerable to small-arms fire, unlike the air-cooled R-2800 radials of its Republic P-47 Thunderbolt stablemates based in England, regularly tasked with ground-strafing missions. Given the overwhelming Allied air superiority, the Luftwaffe put its effort into the development of aircraft of such high performance that they could operate with impunity, but which also made bomber attack much more difficult, merely from the flight velocities they achieved. Foremost among these were the Messerschmitt Me 163B point-defense rocket interceptors, which started their operations with JG 400 near the end of July 1944, and the longer-endurance Messerschmitt Me 262A jet fighter, first flying with the Gruppe-strength Kommando Nowotny unit by the end of September 1944. In action, the Me 163 proved to be more dangerous to the Luftwaffe than to the Allies and was never a serious threat. The Me 262A was a serious threat, but attacks on their airfields neutralized them. The pioneering Junkers Jumo 004 axial-flow jet engines of the Me 262As needed careful nursing by their pilots, and these aircraft were particularly vulnerable during takeoff and landing. Lt. Chuck Yeager of the 357th Fighter Group was one of the first American pilots to shoot down an Me 262, which he caught during its landing approach. On 7 October 1944, Lt. Urban L. Drew of the 361st Fighter Group shot down two Me 262s that were taking off, while on the same day, Lt. Col. Hubert Zemke, who had transferred to the Mustang-equipped 479th Fighter Group, shot down what he thought was a Bf 109, only to have his gun camera film reveal that it may have been an Me 262. On 25 February 1945, Mustangs of the 55th Fighter Group surprised an entire Staffel of Me 262As at takeoff and destroyed six jets. The Mustang also proved useful against the V-1s launched toward London. P-51B/Cs, using 150-octane fuel, were fast enough to catch the V-1 and operated in concert with shorter-range aircraft such as advanced marks of the Supermarine Spitfire and Hawker Tempest. By 8 May 1945, the 8th, 9th, and 15th Air Force's P-51 groups claimed some 4,950 aircraft shot down (about half of all USAAF claims in the European theater, the most claimed by any Allied fighter in air-to-air combat) and 4,131 destroyed on the ground. Losses were about 2,520 aircraft. The 8th Air Force's 4th Fighter Group was the top-scoring fighter group in Europe, with 1,016 enemy aircraft claimed destroyed. This included 550 claimed in aerial combat and 466 on the ground. In air combat, the top-scoring P-51 units (both of which exclusively flew Mustangs) were the 357th Fighter Group of the 8th Air Force with 565 air-to-air combat victories and the 9th Air Force's 354th Fighter Group with 664, which made it one of the top-scoring fighter groups. The top Mustang ace was the USAAF's George Preddy, whose final tally stood at 26.83 victories (a number that includes shared one half- and one third victory credits), 23 of which were scored with the P-51. Preddy was shot down and killed by friendly fire on Christmas Day 1944 during the Battle of the Bulge. In China and the Pacific Theater In early 1945, P-51C, D, and K variants also joined the Chinese Nationalist Air Force. These Mustangs were provided to the 3rd, 4th, and 5th Fighter Groups and used to attack Japanese targets in occupied areas of China. The P-51 became the most capable fighter in China, while the Imperial Japanese Army Air Force used the Nakajima Ki-84 Hayate against it. The P-51 was a relative latecomer to the Pacific theater, due largely to the need for the aircraft in Europe, although the P-38's twin-engined design was considered a safety advantage for long, over-water flights. The first P-51s were deployed in the Far East later in 1944, operating in close-support and escort missions, as well as tactical photoreconnaissance. As the war in Europe wound down, the P-51 became more common. With the capture of Iwo Jima, USAAF P-51 Mustang fighters of the VII Fighter Command were stationed on that island starting in March 1945, being initially tasked with escorting Boeing B-29 Superfortress missions against the Japanese homeland. The command's last major raid of May was a daylight incendiary attack on Yokohama on 29 May conducted by 517 B-29s escorted by 101 P-51s. This force was intercepted by 150 A6M Zero fighters, sparking an intense air battle in which five B-29s were shot down and another 175 damaged. In return, the P-51 pilots claimed 26 "kills" and 23 "probables" for the loss of three fighters. The 454 B-29s that reached Yokohama struck the city's main business district and destroyed of buildings; over 1000 Japanese were killed. Overall, the attacks in May destroyed of buildings, which was equivalent to one-seventh of Japan's total urban area. The minister of home affairs, Iwao Yamazaki, concluded after these raids that Japan's civil defense arrangements were "considered to be futile". On the first day of June, 521 B-29s escorted by 148 P-51s were dispatched in a daylight raid against Osaka. While en route to the city, the Mustangs flew through thick clouds, and 27 of the fighters were destroyed in collisions. Nevertheless, 458 heavy bombers and 27 P-51s reached the city, and the bombardment killed 3,960 Japanese and destroyed of buildings. On 5 June, 473 B-29s struck Kobe by day and destroyed of buildings for the loss of 11 bombers. A force of 409 B-29s attacked Osaka again on 7 June; during this attack, of buildings were burnt out and the Americans did not suffer any losses. Osaka was bombed for the fourth time that month, on 15 June, when 444 B-29s destroyed of the city and another of nearby Amagasaki; 300,000 houses were destroyed in Osaka. This attack marked the end of the first phase of XXI Bomber Command's attack on Japan's cities. During May and June, the bombers had destroyed much of the country's six largest cities, killing between 112,000 and 126,762 people and rendering millions homeless. The widespread destruction and high number of casualties from these raids caused many Japanese to realize that their country's military was no longer able to defend the home islands. American losses were low compared to Japanese casualties; 136 B-29s were downed during the campaign. In Tokyo, Osaka, Nagoya, Yokohama, Kobe, and Kawasaki, "over 126,762 people were killed ... and a million and a half dwellings and over of urban space were destroyed." In Tokyo, Osaka and Nagoya, "the areas leveled (almost ) exceeded the areas destroyed in all German cities by both the American and British air forces (about )." P-51s also conducted a series of independent ground-attack missions against targets in the home islands. The first of these operations took place on 16 April, when 57 P-51s strafed Kanoya Air Field in Kyushu. In operations conducted between 26 April and 22 June, the American fighter pilots claimed the destruction of 64 Japanese aircraft and damage to another 180 on the ground, as well as a further 10 shot down in flight; these claims were lower than the American planners had expected, however, and the raids were considered unsuccessful. USAAF losses were 11 P-51s to enemy action and seven to other causes. Due to the lack of Japanese air opposition to the American bomber raids, VII Fighter Command was solely tasked with ground-attack missions from July. These raids were frequently made against airfields to destroy aircraft being held in reserve to attack the expected Allied invasion fleet. While the P-51 pilots only occasionally encountered Japanese fighters in the air, the airfields were protected by antiaircraft batteries and barrage balloons. By the end of the war, VII Fighter Command had conducted 51 ground-attack raids, of which 41 were considered successful. The fighter pilots claimed to have destroyed or damaged 1,062 aircraft and 254 ships, along with large numbers of buildings and railway rolling stock. American losses were 91 pilots killed and 157 Mustangs destroyed. Medal of Honor recipients Two P-51 pilots received the Medal of Honor during World War II: USAAF Lt Col. James H. Howard of the 356th Fighter Squadron, 354th Fighter Group was awarded the Medal of Honor for his action during a bomber escort mission near Oschersleben, Germany on 11 January 1944, flying P-51B, serial number nicknamed "Ding Hao". Despite being outnumbered, Howard shot down three German planes and continued to defend the bombers even when his guns went out of action and fuel supply became dangerously low. USAAF Maj. William A. Shomo of the 82nd Reconnaissance Squadron, 71st Reconnaissance Group was awarded the Medal of Honor for his action during a mission over Luzon, Philippines on 11 January 1945, flying an F-6D, the armed photo reconnaissance variant of the P-51, serial number nicknamed "Snooks the 5th". On that mission, Shomo shot down seven Japanese planes and became an "ace in a day". Pilot observations Chief Naval Test Pilot and C.O. Captured Enemy Aircraft Flight Capt. Eric Brown, RN, tested the Mustang at RAE Farnborough in March 1944 and noted: The U.S. Air Forces, Flight Test Engineering, assessed the Mustang B on 24 April 1944 thus: Kurt Bühligen, the third-highest scoring German fighter pilot of World War II's Western Front (with 112 confirmed victories, three against Mustangs), later stated: German fighter ace Heinz Bär said that the P-51: After World War II In the aftermath of World War II, the USAAF consolidated much of its wartime combat force and selected the P-51 as a "standard" piston-engined fighter, while other types, such as the P-38 and P-47, were withdrawn or given substantially reduced roles. As the more advanced (P-80 and P-84) jet fighters were introduced, the P-51 was also relegated to secondary duties. In 1947, the newly formed USAF Strategic Air Command employed Mustangs alongside F-6 Mustangs and F-82 Twin Mustangs, due to their range capabilities. In 1948, the designation P-51 (P for pursuit) was changed to F-51 (F for fighter) and the existing F designator for photographic reconnaissance aircraft was dropped because of a new designation scheme throughout the USAF. Aircraft still in service in the USAF or Air National Guard (ANG) when the system was changed included: F-51B, F-51D, F-51K, RF-51D (formerly F-6D), RF-51K (formerly F-6K) and TRF-51D (two-seat trainer conversions of F-6Ds). They remained in service from 1946 through 1951. By 1950, although Mustangs continued in service with the USAF after the war, the majority of the USAF's Mustangs had become surplus to requirements and placed in storage, while some were transferred to the Air Force Reserve and the ANG. From the start of the Korean War, the Mustang once again proved useful. A "substantial number" of stored or in-service F-51Ds were shipped, via aircraft carriers, to the combat zone, and were used by the USAF, the South African Air Force, and the Republic of Korea Air Force (ROKAF). The F-51 was used for ground attack, fitted with rockets and bombs, and photo reconnaissance, rather than being as interceptors or "pure" fighters, where it was already surpassed by early jets. After the first North Korean invasion, USAF units were forced to fly from bases in Japan and the F-51Ds, with their long range and endurance, could attack targets in Korea that short-ranged F-80 jets could not. Because of the vulnerable liquid cooling system, however, the F-51s sustained heavy losses to ground fire. Due to its lighter structure and a shortage of spare parts, the newer, faster F-51H was not used in Korea. On August 5, 1950, Major Louis J. Sebille of the 67th Fighter-Bomber Squadron attacked a North Korean armored column advancing on United Nations military units during the Battle of Pusan Perimeter. Though his aircraft was heavily damaged and he was wounded during the first pass on the column, he turned his F-51 around and deliberately crashed into the convoy at the cost of his life, and was posthumously awarded the Medal of Honor. Mustangs continued flying with USAF and ROKAF fighter-bomber units on close support and interdiction missions in Korea until 1953 when they were largely replaced as fighter-bombers by USAF F-84s and by United States Navy (USN) Grumman F9F Panthers. Other air forces and units using the Mustang included the Royal Australian Air Force's 77 Squadron, which flew Australian-built Mustangs as part of British Commonwealth Forces Korea. The Mustangs were replaced by Gloster Meteor F8s in 1951. The South African Air Force's 2 Squadron used U.S.-built Mustangs as part of the U.S. 18th Fighter Bomber Wing and had suffered heavy losses by 1953, after which 2 Squadron converted to the F-86 Sabre. F-51s flew in the Air Force Reserve and ANG throughout the 1950s. The last American USAF Mustang was F-51D-30-NA AF serial no. , which was finally withdrawn from service with the West Virginia Air National Guard's 167th Fighter Interceptor Squadron in January 1957 and retired to what was then called the Air Force Central Museum, although it was briefly reactivated to fly at the 50th anniversary of the Air Force Aerial Firepower Demonstration at the Air Proving Ground, Eglin AFB, Florida, on 6 May 1957. This aircraft, painted as P-51D-15-NA serial no. , is on display at the National Museum of the United States Air Force, Wright-Patterson AFB, in Dayton, Ohio. The final withdrawal of the Mustang from USAF dumped hundreds of P-51s onto the civilian market. The rights to the Mustang design were purchased from North American by the Cavalier Aircraft Corporation, which attempted to market the surplus Mustang aircraft in the U.S. and overseas. In 1967 and again in 1972, the USAF procured batches of remanufactured Mustangs from Cavalier, most of them destined for air forces in South America and Asia that were participating in the Military Assistance Program (MAP). These aircraft were remanufactured from existing original F-51D airframes fitted with new V-1650-7 engines, a new radio, tall F-51H-type vertical tails, and a stronger wing that could carry six machine guns and a total of eight underwing hardpoints. Two bombs and six rockets could be carried. They all had an original F-51D-type canopy but carried a second seat for an observer behind the pilot. One additional Mustang was a two-seat, dual-control TF-51D (67-14866) with an enlarged canopy and only four wing guns. Although these remanufactured Mustangs were intended for sale to South American and Asian nations through the MAP, they were delivered to the USAF with full USAF markings. They were, however, allocated new serial numbers (, and ). The last U.S. military use of the F-51 was in 1968 when the U. S. Army employed a vintage F-51D (44-72990) as a chase aircraft for the Lockheed YAH-56 Cheyenne armed helicopter project. This aircraft was so successful that the Army ordered two F-51Ds from Cavalier in 1968 for use at Fort Rucker as chase planes. They were assigned the serials and . These F-51s had wingtip fuel tanks and were unarmed. Following the end of the Cheyenne program, these two chase aircraft were used for other projects. One of them (68-15795) was fitted with a 106 mm recoilless rifle for evaluation of the weapon's value in attacking fortified ground targets. Cavalier Mustang 68-15796 survives at the Air Force Armament Museum, Eglin AFB, Florida, displayed indoors in World War II markings. The F-51 was adopted by many foreign air forces and continued to be an effective fighter into the mid-1980s with smaller air arms. The last Mustang ever downed in battle occurred during Operation Power Pack in the Dominican Republic in 1965, with the last aircraft finally being retired by the Dominican Air Force in 1984. Service with other air forces After World War II, the P-51 Mustang served in the air arms of more than 25 nations. During the war, a Mustang cost about $51,000, while many hundreds were sold postwar for the nominal price of one dollar to signatories of the Inter-American Treaty of Reciprocal Assistance, ratified in Rio de Janeiro in 1947. These countries used the P-51 Mustang: Australia In November 1944, 3 Squadron RAAF became the first Royal Australian Air Force unit to use Mustangs. At the time of its conversion from the P-40 to the Mustang, the squadron was based in Italy with the RAF's First Tactical Air Force. 3 Squadron was renumbered 4 Squadron after returning to Australia from Italy, and converted to P-51Ds. Several other Australian or Pacific-based squadrons converted to either CAC-built Mustangs or to imported P-51Ks from July 1945, having been equipped with P-40s or Boomerangs for wartime service; these units were: 76, 77, 82, 83, 84 and 86 squadrons. Only 17 Mustangs reached the RAAF's First Tactical Air Force front-line squadrons by the time World War II ended in August 1945. 76, 77 and 82 squadrons were formed into 81 Fighter Wing of the British Commonwealth Air Force, which was part of the British Commonwealth Occupation Force stationed in Japan from February 1946. 77 Squadron used its P-51s extensively during the first months of the Korean War, before converting to Gloster Meteor jets. Five reserve units from the Citizen Air Force also operated Mustangs. 21 "City of Melbourne" Squadron, based in the state of Victoria; 22 "City of Sydney" Squadron, based in New South Wales; 23 "City of Brisbane" Squadron, based in Queensland; 24 "City of Adelaide" Squadron, based in South Australia; and 25 "City of Perth" Squadron, based in Western Australia; all of these units were equipped with CAC Mustangs, rather than P-51D or Ks. The last Mustangs were retired from these units in 1960 when CAF units adopted a nonflying role. Bolivia Nine Cavalier F-51D (including the two TF-51s) were given to Bolivia, under a program called Peace Condor. Canada Canada had five squadrons equipped with Mustangs during World War II. RCAF 400, 414, and 430 squadrons flew Mustang Mk Is (1942–1944) and 441 and 442 squadrons flew Mustang Mk IIIs and IVAs in 1945. Postwar, a total of 150 Mustang P-51Ds were purchased and served in two regular (416 "Lynx" and 417 "City of Windsor") and six auxiliary fighter squadrons (402 "City of Winnipeg", 403 "City of Calgary", 420 "City of London", 424 "City of Hamilton", 442 "City of Vancouver" and 443 "City of New Westminster"). The Mustangs were declared obsolete in 1956, but a number of special-duty versions served on into the early 1960s. Republic of China The Chinese Nationalist Air Force obtained the P-51 during the late Sino-Japanese War to fight against the Japanese. After the war, Chiang Kai-shek's Nationalist government used the planes against insurgent Communist forces. The Nationalists retreated to Taiwan in 1949. Pilots supporting Chiang brought most of the Mustangs with them, where the aircraft became part of the island's defense arsenal. People's Republic of China The Communist Chinese captured 39 P-51s from the Nationalists while they were retreating to Taiwan. In August 1949, the People's Liberation Army Air Force formed its first P-51 squadron at Beijing Nanyuan Airport and were tasked of the defending Beijing's airspace from Nationalist Air Force aircraft. On 1 October 1949, when Mao Zedong proclaimed the founding of the People's Republic of China, nine P-51s conducted a fly-past during the military parade in Beijing. By 1950, when Soviet Union began supplying modern military equipment to China, surviving P-51s were relegated to PLAAF's aviation school and 13 P-51s were modified as two-seat trainers. By September 1953, most P-51s were retired from service and only eight P-51s remained in service to teach Ilyushin Il-10 pilots on how to taxi aircraft. Costa Rica The Costa Rican Air Force flew four P-51Ds from 1955 to 1964. Cuba In November 1958, three US-registered civilian P-51D Mustangs were illegally flown separately from Miami to Cuba, on delivery to the rebel forces of the 26th of July Movement, then headed by Fidel Castro during the Cuban Revolution. One of the Mustangs was damaged during delivery and none of them were used operationally. After the success of the revolution in January 1959, with other rebel aircraft plus those of the existing Cuban government forces, they were adopted into the Fuerza Aérea Revolucionaria. Due to increasing U.S. restrictions and lack of spares and maintenance experience, they never achieved operational status. At the time of the Bay of Pigs invasion, the two intact Mustangs were already effectively grounded at Campo Columbia and at Santiago. After the failed invasion, they were placed on display with other symbols of "revolutionary struggle" and one remains on display at the Museo del Aire. Dominican Republic The Dominican Republic was the largest Latin American air force to employ the P-51D, with six aircraft acquired in 1948, 44 ex-Swedish F-51Ds purchased in 1948, and a further Mustang obtained from an unknown source. It was the last nation to have any Mustangs in service, with some remaining in use as late as 1984. Nine of the final 10 aircraft were sold back to American collectors in 1988. El Salvador The Salvadoran Air Force (Fuerza Aérea Salvadoreña or FAS) purchased five Cavalier Mustang IIs (and one dual-control Cavalier TF-51) that featured wingtip fuel tanks to increase combat range and up-rated Merlin engines. Seven P-51D Mustangs were also in service. They were used during the 1969 Football War against Honduras, the last time the P-51 was used in combat. One of them, FAS-404, was shot down by a Vought F4U-5 Corsair flown by Captain Fernando Soto in the last aerial combat between piston-engined fighters in the world. France In late 1944, the first French unit began its transition to reconnaissance Mustangs. In January 1945, the Tactical Reconnaissance Squadron 2/33 of the French Air Force took their F-6Cs and F-6Ds over Germany on photographic mapping missions. The Mustangs remained in service until the early 1950s, when they were replaced by jet fighters. Germany Several P-51s were captured by the Luftwaffe as Beuteflugzeug ("captured aircraft") following crash landings. These aircraft were subsequently repaired and test-flown by the Zirkus Rosarius, or Rosarius Staffel, the official Erprobungskommando of the Luftwaffe High Command, for combat evaluation at Göttingen. The aircraft were repainted with German markings and bright yellow noses, tails, and bellies for identification. A number of P-51B/P-51Cs – including examples marked with Luftwaffe Geschwaderkennung codes T9+CK, T9+FK, T9+HK, and T9+PK (with the "T9" prefix not known to be officially assigned to any existing Luftwaffe formation from their own records, outside of the photos of Zirkus Rosarius–flown aircraft)—with a total of three captured P-51Ds were also flown by the unit. Some of these P-51s were found by Allied forces at the end of the war; others crashed during testing. The Mustang is also listed in the appendix to the novel KG 200 as having been flown by the German secret operations unit KG 200, which tested, evaluated, and sometimes clandestinely operated captured enemy aircraft during World War II. Guatemala The Guatemalan Air Force had 30 P-51D Mustangs in service from 1954 to the early 1970s. Haiti Haiti had four P-51D Mustangs when President Paul Eugène Magloire was in power from 1950 to 1956, with the last retired in 1973–1974 and sold for spares to the Dominican Republic. Indonesia Indonesia acquired 26 P-51D/Ks from the departing Netherlands East Indies Air Force in 1949–1950 and later received 35 P-51Ds from the United States in 1960–1961. The Mustangs were used against numerous rebellions during the 1950s, such as the CIA-backed Permesta rebels in 1958–1961. During this period, the Mustang scored the first and (as of 2022) the only aerial victory of the Indonesian Air Force, when on 18 May 1958, a P-51D Mustang piloted by Capt. Ignatius Dewanto shot down a Permesta's Revolutionary Air Force B-26 Invader piloted by Allen Lawrence Pope near Ambon. They were also used against Commonwealth (RAF, RAAF, and RNZAF) forces during the Indonesia–Malaysia confrontation in the early 1960s. Indonesia received a shipment of five or seven Cavalier II Mustangs and one TF-51D (without tip tanks) delivered in 1972–1973 as part of "Peace Pony" program under the Mutual Defense Assistance Act. The last time Mustangs were deployed for military purposes was during the "Wibawa V" exercise at Mount Lawu, Magetan in February 1975. The Indonesian Mustangs were also used for filming Janur Kuning, which was released in 1980. The Mustangs were replaced in 1976. Israel A few P-51 Mustangs were illegally bought by Israel in 1948, crated, and smuggled into the country as agricultural equipment for use in the 1947–1949 Palestine war, serving alongside upwards of 23 Avia S-199 fighters (Czech-built Messerschmitt Bf 109Gs) in Israeli service, with the Mustangs quickly establishing themselves as the best fighter in the Israeli inventory. Further aircraft were bought from Sweden and were replaced by jets at the end of the 1950s, but not before the type was used in the Suez Crisis, at the opening of Operation Kadesh. In conjunction with a surprise parachute drop at the Mitla Pass, four P-51s were specially detailed to cut telephone and telegraph wires using their wings in extreme low level runs, which resulted in major interruptions to Egyptian communications. Italy Italy was a postwar operator of P-51Ds; deliveries were slowed by the Korean War, but between September 1947 and January 1951, by MDAP count, 173 examples were delivered. They were used in all the AMI fighter units: 2, 3, 4, 5, 6 and 51 Stormo (wing), plus some employed in schools and experimental units. Considered a "glamorous" fighter, P-51s were even used as personal aircraft by several Italian commanders. Some restrictions were placed on its use due to unfavorable flying characteristics. Handling had to be done with much care when fuel tanks were fully used, and several aerobatic maneuvers were forbidden. Overall, the P-51D was highly rated even compared to the other primary postwar fighter in Italian service, the Supermarine Spitfire, partly because these P-51Ds were in very good condition in contrast to all other Allied fighters supplied to Italy. Phasing out of the Mustang began in summer 1958. Japan The P-51C-11-NT Evalina, marked as "278" (former USAAF serial: ) and flown by 26th FS, 51st FG, was hit by gunfire on 16 January 1945 and belly-landed on Suchon Airfield in China, which was held by the Japanese. The Japanese repaired the aircraft, roughly applied Hinomaru roundels and flew the aircraft to the Fussa evaluation center (now Yokota Air Base) in Japan. Netherlands The Royal Netherlands East Indies Army Air Force received 40 P-51Ds and flew them in the course of the Indonesian National Revolution, particularly during the two Dutch police actions: Operation Product in 1947 and Operation Kraai in 1948–1949. When the conflict was over, Indonesia received 26 of these Mustangs. New Zealand New Zealand ordered 370 P-51 Mustangs to supplement its F4U Corsairs in the Pacific Ocean Areas theater. Scheduled deliveries were for an initial batch of 30 P-51Ds, followed by 137 more P-51Ds and 203 P-51Ms. The original 30 were being shipped as the war ended in August 1945; these were stored in their packing cases, and the order for the additional Mustangs was canceled. In 1951, the stored Mustangs entered service in 1 (Auckland), 2 (Wellington), 3 (Canterbury), and 4 (Otago) squadrons of the Territorial Air Force (TAF). The Mustangs remained in service until they were prematurely retired in August 1955 following a series of problems with undercarriage and coolant-system corrosion problems. Four Mustangs served on as target tugs until the TAF was disbanded in 1957. RNZAF pilots in the Royal Air Force also flew the P-51 and at least one New Zealand pilot scored victories over Europe while on loan to a USAAF P-51 squadron. Nicaragua The Nicaraguan National Guard purchased 26 P-51D Mustangs from Sweden in 1954 and later received 30 P-51D Mustangs from the U.S. together with two TF-51 models from MAP after 1954. All aircraft of this type were retired from service by 1964. Philippines The Philippines acquired 103 P-51D Mustangs after World War II, operated by the 6th "Cobras", 7th "Bulldogs" and 8th "Scorpions" tactical fighter squadrons of the 5th Fighter Wing. These became the backbone of the postwar Philippine Army Air Corps and Philippine Air Force, and were used extensively during the Huk campaign, fighting against communist insurgents, as well as the suppression of Moro rebels led by Hadji Kamlon in southern Philippines until 1955. The Mustangs were also the first aircraft of the Philippine air demonstration team, which was formed in 1953 and given the name the "Blue Diamonds" the following year. The Mustangs were replaced by 56 F-86 Sabres in the late 1950s, but some were still in service for COIN roles up to the early 1980s. Poland During World War II, five Polish Air Force in Great Britain squadrons used Mustangs. The first Polish unit equipped (7 June 1942) with Mustang Mk Is was "B" Flight of 309 "Ziemi Czerwieńskiej" Squadron (an Army Co-Operation Command unit), followed by "A" Flight in March 1943. Subsequently, 309 Squadron was redesignated a fighter/reconnaissance unit and became part of Fighter Command. On 13 March 1944, 316 "Warszawski" Squadron received their first Mustang Mk IIIs; rearming of the unit was completed by the end of April. By 26 March 1944, 306 "Toruński" Sqn and 315 "Dębliński" Sqn received Mustangs Mk IIIs (the whole operation took 12 days). On 20 October 1944, Mustang Mk Is in 309 Squadron were replaced by Mk IIIs. On 11 December 1944, the unit was again renamed, becoming 309 Dywizjon Myśliwski "Ziemi Czerwieńskiej" or 309 "Land of Czerwien" Polish Fighter Squadron. In 1945, 303 "Kościuszko" Sqn received 20 Mustangs Mk IV/Mk IVA replacements. Postwar, between 6 December 1946 and 6 January 1947, all five Polish squadrons equipped with Mustangs were disbanded. Poland returned about 80 Mustang Mk IIIs and 20 Mustangs Mk IV/IVAs to the RAF, which transferred them to the U.S. government. Somalia The Somalian Air Force operated eight P-51Ds in post-World War II service. South Africa No.5 Squadron South African Air Force operated a number of Mustang Mk IIIs (P-51B/C) and Mk IVs (P-51D/K) in Italy during World War II, beginning in September 1944, when the squadron converted to the Mustang Mk III from Kittyhawks. The Mk IV and Mk IVA came into SA service in March 1945. These aircraft were generally camouflaged in the British style, having been drawn from RAF stocks; all carried RAF serial numbers and were struck off charge and scrapped in October 1945. In 1950, 2 Squadron SAAF was supplied with F-51D Mustangs by the United States for Korean War service. The type performed well in South African hands before being replaced by the F-86 Sabre in 1952 and 1953. South Korea Within a month of the outbreak of the Korean War, 10 F-51D Mustangs were provided to the badly depleted Republic of Korea Air Force as a part of the Bout One Project. They were flown by both South Korean airmen, several of whom were veterans of the Imperial Japanese Army and Navy air services during World War II, as well as by U.S. advisers led by Major Dean Hess. Later, more were provided both from U.S. and from South African stocks, as the latter were converting to F-86 Sabres. They formed the backbone of the South Korean Air Force until they were replaced by Sabres. It also served with the ROKAF Black Eagles aerobatic team, until retired in 1954. Sweden Sweden's Flygvapnet first recuperated four of the P-51s (two P-51Bs and two early P-51Ds) that had been diverted to Sweden during missions over Europe. In February 1945, Sweden purchased 50 P-51Ds designated J 26, which were delivered by American pilots in April and assigned to the Uppland Wing (F 16) at Uppsala as interceptors. In early 1946, the Jämtland Wing (F 4) at Östersund was equipped with a second batch of 90 P-51Ds. A final batch of 21 Mustangs was purchased in 1948. In all, 161 J 26s served in the Swedish Air Force during the late 1940s. About 12 were modified for photo reconnaissance and redesignated S 26. Some of these aircraft participated in the secret Swedish mapping of new Soviet military installations at the Baltic coast in 1946–47 (Operation Falun), an endeavor that entailed many intentional violations of Soviet airspace. However, the Mustang could outdive any Soviet fighter of that era, so no S 26s were lost in these missions. The J 26s were replaced by De Havilland Vampires around 1950. The S 26s were replaced by S 29Cs in the early 1950s. Switzerland The Swiss Air Force operated a few USAAF P-51s that had been impounded by Swiss authorities during World War II after the pilots were forced to land in neutral Switzerland. After the war, Switzerland also bought 130 P-51s for $4,000 each. They served until 1958. Soviet Union The Soviet Union received at least 10 early-model ex-RAF Mustang Mk Is and tested them, but found them to "under-perform" compared to contemporary USSR fighters, relegating them to training units. Later Lend-Lease deliveries of the P-51B/C and D series, along with other Mustangs abandoned in Russia after the famous "shuttle missions", were repaired and used by the Soviet Air Force, but not in front-line service. Uruguay The Uruguayan Air Force used 25 P-51D Mustangs from 1950 to 1960; some were subsequently sold to Bolivia. P-51s and civil aviation Many P-51s were sold as surplus after the war, often for as little as $1,500. Some were sold to former wartime fliers or other aficionados for personal use, while others were modified for air racing. One of the most significant Mustangs involved in air racing was serial number , a surplus P-51C-10-NT purchased by film stunt pilot Paul Mantz. He modified the wings, sealing them to create a giant fuel tank in each one; these "wet wings" reduced the need for fuel stops or drag-inducing drop tanks. Named Blaze of Noon after the film Blaze of Noon, the aircraft won the 1946 and 1947 Bendix Air Races, took second in the 1948 Bendix, and placed third in the 1949 Bendix. Mantz also set a U.S. coast-to-coast record in 1947. He sold the Mustang to Charles F. Blair Jr (future husband of Maureen O'Hara), who renamed it Excalibur III and used it to set a New York-to-London (about ) record in 1951: 7 hr 48 min from takeoff at Idlewild to overhead London Airport. Later that year, Blair flew from Norway to Fairbanks, Alaska, via the North Pole (about ), proving that navigation via sun sights was possible over the magnetic North Pole region. For this feat, he was awarded the Harmon Trophy and the Air Force was forced to change its thoughts on a possible Soviet air strike from the north. This Mustang now sits in the National Air and Space Museum's Steven F. Udvar-Hazy Center. In 1958, the RCAF retired its 78 remaining Mustangs. RCAF pilot Lynn Garrison ferried them from their various storage locations to Canastota, New York, where the American buyers were based. Garrison flew each of the surviving aircraft at least once. These aircraft make up a large percentage of the aircraft presently flying worldwide. The most prominent firm to convert Mustangs to civilian use was Trans-Florida Aviation, later renamed Cavalier Aircraft Corporation, which produced the Cavalier Mustang. Modifications included a taller tailfin and wingtip tanks. A number of conversions included a Cavalier Mustang specialty: a "tight" second seat added in the space formerly occupied by the military radio and fuselage fuel tank. In the late 1960s and early 1970s, when the United States Department of Defense wished to supply aircraft to South American countries and later Indonesia for close air support and counterinsurgency, it paid Cavalier to return some of their civilian conversions back to updated military specifications. In the 21st century, a P-51 can command a price of more than $1 million, even for only partially restored aircraft. There were 204 privately owned P-51s in the U.S. on the FAA registry in 2011, most of which are still flying, often associated with organizations such as the Commemorative Air Force (formerly the Confederate Air Force). In May 2013, Doug Matthews set an altitude record of in a P-51 named The Rebel for piston-powered aircraft weighing . Flying from a grass runway at Florida's Indiantown airport and over Lake Okeechobee, Matthews set world records for time to reach altitudes of , 18 minutes and , 31 minutes. He set a level-flight altitude record of in level flight and an absolute altitude record of , breaking the previous record of set in 1954. Incidents On 9 June 1973, William Penn Patrick (43) a certified pilot and his passenger, Christian Hagert, died when Patrick's P-51 Mustang crashed in Lakeport, California. On 1 July 1990 at the National Capital Air Show (Ottawa, Ontario, Canada), Harry E. Tope was killed when his P-51 Mustang crashed. On 16 September 2011 The Galloping Ghost, a modified P-51 piloted by Jimmy Leeward of Ocala, Florida, crashed during an air race in Reno, Nevada. Leeward and at least nine people on the ground were killed when the racer suddenly crashed near the edge of the grandstand. Variants Over 20 variants of the P-51 Mustang were produced from 1940 to after the war. Production Except for the small numbers assembled or produced in Australia, all Mustangs were built by North American initially at Inglewood, California, but then additionally in Dallas, Texas. Accidents and incidents Surviving aircraft Specifications (P-51D Mustang) Notable appearances in media Red Tail Reborn (2007) is the story behind the restoration of a flying memorial aircraft. Scale replicas As indicative of the iconic nature of the P-51, manufacturers within the hobby industry have created scale plastic model kits of the P-51 Mustang, with varying degrees of detail and skill levels. The aircraft have also been the subject of numerous scale flying replicas. Aside from the popular model aircraft, several kitplane manufacturers offer ½, ⅔, and ¾-scale replicas capable of comfortably seating one (or even two) and offering high performance combined with more forgiving flight characteristics. Such aircraft include the Titan T-51 Mustang, W.A.R. P-51 Mustang, Linn Mini Mustang, Jurca Gnatsum, Thunder Mustang, Stewart S-51D Mustang, Loehle 5151 Mustang and ScaleWings SW51 Mustang.	Technology	Specific aircraft	null
24714	https://en.wikipedia.org/wiki/Precession	Precession	Precession is a change in the orientation of the rotational axis of a rotating body. In an appropriate reference frame it can be defined as a change in the first Euler angle, whereas the third Euler angle defines the rotation itself. In other words, if the axis of rotation of a body is itself rotating about a second axis, that body is said to be precessing about the second axis. A motion in which the second Euler angle changes is called nutation. In physics, there are two types of precession: torque-free and torque-induced. In astronomy, precession refers to any of several slow changes in an astronomical body's rotational or orbital parameters. An important example is the steady change in the orientation of the axis of rotation of the Earth, known as the precession of the equinoxes. Torque-free or torque neglected Torque-free precession implies that no external moment (torque) is applied to the body. In torque-free precession, the angular momentum is a constant, but the angular velocity vector changes orientation with time. What makes this possible is a time-varying moment of inertia, or more precisely, a time-varying inertia matrix. The inertia matrix is composed of the moments of inertia of a body calculated with respect to separate coordinate axes (e.g. , , ). If an object is asymmetric about its principal axis of rotation, the moment of inertia with respect to each coordinate direction will change with time, while preserving angular momentum. The result is that the component of the angular velocities of the body about each axis will vary inversely with each axis' moment of inertia. The torque-free precession rate of an object with an axis of symmetry, such as a disk, spinning about an axis not aligned with that axis of symmetry can be calculated as follows: where is the precession rate, is the spin rate about the axis of symmetry, is the moment of inertia about the axis of symmetry, is moment of inertia about either of the other two equal perpendicular principal axes, and is the angle between the moment of inertia direction and the symmetry axis. When an object is not perfectly rigid, inelastic dissipation will tend to damp torque-free precession, and the rotation axis will align itself with one of the inertia axes of the body. For a generic solid object without any axis of symmetry, the evolution of the object's orientation, represented (for example) by a rotation matrix that transforms internal to external coordinates, may be numerically simulated. Given the object's fixed internal moment of inertia tensor and fixed external angular momentum , the instantaneous angular velocity is Precession occurs by repeatedly recalculating and applying a small rotation vector for the short time ; e.g.: for the skew-symmetric matrix . The errors induced by finite time steps tend to increase the rotational kinetic energy: this unphysical tendency can be counteracted by repeatedly applying a small rotation vector perpendicular to both and , noting that Torque-induced Torque-induced precession (gyroscopic precession) is the phenomenon in which the axis of a spinning object (e.g., a gyroscope) describes a cone in space when an external torque is applied to it. The phenomenon is commonly seen in a spinning toy top, but all rotating objects can undergo precession. If the speed of the rotation and the magnitude of the external torque are constant, the spin axis will move at right angles to the direction that would intuitively result from the external torque. In the case of a toy top, its weight is acting downwards from its center of mass and the normal force (reaction) of the ground is pushing up on it at the point of contact with the support. These two opposite forces produce a torque which causes the top to precess. The device depicted on the right is gimbal mounted. From inside to outside there are three axes of rotation: the hub of the wheel, the gimbal axis, and the vertical pivot. To distinguish between the two horizontal axes, rotation around the wheel hub will be called spinning, and rotation around the gimbal axis will be called pitching. Rotation around the vertical pivot axis is called rotation. First, imagine that the entire device is rotating around the (vertical) pivot axis. Then, spinning of the wheel (around the wheelhub) is added. Imagine the gimbal axis to be locked, so that the wheel cannot pitch. The gimbal axis has sensors, that measure whether there is a torque around the gimbal axis. In the picture, a section of the wheel has been named . At the depicted moment in time, section is at the perimeter of the rotating motion around the (vertical) pivot axis. Section , therefore, has a lot of angular rotating velocity with respect to the rotation around the pivot axis, and as is forced closer to the pivot axis of the rotation (by the wheel spinning further), because of the Coriolis effect, with respect to the vertical pivot axis, tends to move in the direction of the top-left arrow in the diagram (shown at 45°) in the direction of rotation around the pivot axis. Section of the wheel is moving away from the pivot axis, and so a force (again, a Coriolis force) acts in the same direction as in the case of . Note that both arrows point in the same direction. The same reasoning applies for the bottom half of the wheel, but there the arrows point in the opposite direction to that of the top arrows. Combined over the entire wheel, there is a torque around the gimbal axis when some spinning is added to rotation around a vertical axis. It is important to note that the torque around the gimbal axis arises without any delay; the response is instantaneous. In the discussion above, the setup was kept unchanging by preventing pitching around the gimbal axis. In the case of a spinning toy top, when the spinning top starts tilting, gravity exerts a torque. However, instead of rolling over, the spinning top just pitches a little. This pitching motion reorients the spinning top with respect to the torque that is being exerted. The result is that the torque exerted by gravity – via the pitching motion – elicits gyroscopic precession (which in turn yields a counter torque against the gravity torque) rather than causing the spinning top to fall to its side. Precession or gyroscopic considerations have an effect on bicycle performance at high speed. Precession is also the mechanism behind gyrocompasses. Classical (Newtonian) Precession is the change of angular velocity and angular momentum produced by a torque. The general equation that relates the torque to the rate of change of angular momentum is: where and are the torque and angular momentum vectors respectively. Due to the way the torque vectors are defined, it is a vector that is perpendicular to the plane of the forces that create it. Thus it may be seen that the angular momentum vector will change perpendicular to those forces. Depending on how the forces are created, they will often rotate with the angular momentum vector, and then circular precession is created. Under these circumstances the angular velocity of precession is given by: where is the moment of inertia, is the angular velocity of spin about the spin axis, is the mass, is the acceleration due to gravity, is the angle between the spin axis and the axis of precession and is the distance between the center of mass and the pivot. The torque vector originates at the center of mass. Using , we find that the period of precession is given by: Where is the moment of inertia, is the period of spin about the spin axis, and is the torque. In general, the problem is more complicated than this, however. Relativistic (Einsteinian) The special and general theories of relativity give three types of corrections to the Newtonian precession, of a gyroscope near a large mass such as Earth, described above. They are: Thomas precession, a special-relativistic correction accounting for an object (such as a gyroscope) being accelerated along a curved path. de Sitter precession, a general-relativistic correction accounting for the Schwarzschild metric of curved space near a large non-rotating mass. Lense–Thirring precession, a general-relativistic correction accounting for the frame dragging by the Kerr metric of curved space near a large rotating mass. The Schwarzschild geodesics (sometimes Schwarzschild precession) is used in the prediction of the anomalous perihelion precession of the planets, most notably for the accurate prediction of the apsidal precession of Mercury. Astronomy In astronomy, precession refers to any of several gravity-induced, slow and continuous changes in an astronomical body's rotational axis or orbital path. Precession of the equinoxes, perihelion precession, changes in the tilt of Earth's axis to its orbit, and the eccentricity of its orbit over tens of thousands of years are all important parts of the astronomical theory of ice ages. (See Milankovitch cycles.) Axial precession (precession of the equinoxes) Axial precession is the movement of the rotational axis of an astronomical body, whereby the axis slowly traces out a cone. In the case of Earth, this type of precession is also known as the precession of the equinoxes, lunisolar precession, or precession of the equator. Earth goes through one such complete precessional cycle in a period of approximately 26,000 years or 1° every 72 years, during which the positions of stars will slowly change in both equatorial coordinates and ecliptic longitude. Over this cycle, Earth's north axial pole moves from where it is now, within 1° of Polaris, in a circle around the ecliptic pole, with an angular radius of about 23.5°. The ancient Greek astronomer Hipparchus (c. 190–120 BC) is generally accepted to be the earliest known astronomer to recognize and assess the precession of the equinoxes at about 1° per century (which is not far from the actual value for antiquity, 1.38°), although there is some minor dispute about whether he was. In ancient China, the Jin-dynasty scholar-official Yu Xi ( 307–345 AD) made a similar discovery centuries later, noting that the position of the Sun during the winter solstice had drifted roughly one degree over the course of fifty years relative to the position of the stars. The precession of Earth's axis was later explained by Newtonian physics. Being an oblate spheroid, Earth has a non-spherical shape, bulging outward at the equator. The gravitational tidal forces of the Moon and Sun apply torque to the equator, attempting to pull the equatorial bulge into the plane of the ecliptic, but instead causing it to precess. The torque exerted by the planets, particularly Jupiter, also plays a role. Apsidal precession The orbits of planets around the Sun do not really follow an identical ellipse each time, but actually trace out a flower-petal shape because the major axis of each planet's elliptical orbit also precesses within its orbital plane, partly in response to perturbations in the form of the changing gravitational forces exerted by other planets. This is called perihelion precession or apsidal precession. In the adjunct image, Earth's apsidal precession is illustrated. As the Earth travels around the Sun, its elliptical orbit rotates gradually over time. The eccentricity of its ellipse and the precession rate of its orbit are exaggerated for visualization. Most orbits in the Solar System have a much smaller eccentricity and precess at a much slower rate, making them nearly circular and nearly stationary. Discrepancies between the observed perihelion precession rate of the planet Mercury and that predicted by classical mechanics were prominent among the forms of experimental evidence leading to the acceptance of Einstein's Theory of Relativity (in particular, his General Theory of Relativity), which accurately predicted the anomalies. Deviating from Newton's law, Einstein's theory of gravitation predicts an extra term of , which accurately gives the observed excess turning rate of 43 arcseconds every 100 years. Nodal precession Orbital nodes also precess over time.	Physical sciences	Classical mechanics	Physics
24718	https://en.wikipedia.org/wiki/Ring%20system	Ring system	A ring system is a disc or torus orbiting an astronomical object that is composed of solid material such as dust, meteoroids, planetoids, moonlets, or stellar objects. Ring systems are best known as planetary rings, common components of satellite systems around giant planets such as the rings of Saturn, or circumplanetary disks. But they can also be galactic rings and circumstellar discs, belts of planetoids, such as the asteroid belt or Kuiper belt, or rings of interplanetary dust, such as around the Sun at distances of Mercury, Venus, and Earth, in mean motion resonance with these planets. Evidence suggests that ring systems may also be found around other types of astronomical objects, including moons and brown dwarfs. In the Solar System, all four giant planets (Jupiter, Saturn, Uranus, and Neptune) have ring systems. Ring systems around minor planets have also been discovered via occultations. Some studies even theorize that the Earth may have had a ring system during the mid-late Ordovician period. Formation There are three ways that thicker planetary rings have been proposed to have formed: from material originating from the protoplanetary disk that was within the Roche limit of the planet and thus could not coalesce to form moons, from the debris of a moon that was disrupted by a large impact, or from the debris of a moon that was disrupted by tidal stresses when it passed within the planet's Roche limit. Most rings were thought to be unstable and to dissipate over the course of tens or hundreds of millions of years, but it now appears that Saturn's rings might be quite old, dating to the early days of the Solar System. Fainter planetary rings can form as a result of meteoroid impacts with moons orbiting around the planet or, in the case of Saturn's E-ring, the ejecta of cryovolcanic material. Ring systems may form around centaurs when they are tidally disrupted in a close encounter (within 0.4 to 0.8 times the Roche limit) with a giant planet. For a differentiated body approaching a giant planet at an initial relative velocity of 3−6 km/s with an initial rotational period of 8 hours, a ring mass of 0.1%−10% of the centaur's mass is predicted. Ring formation from an undifferentiated body is less likely. The rings would be composed mostly or entirely of material from the parent body's icy mantle. After forming, the ring would spread laterally, leading to satellite formation from whatever portion of it spreads beyond the centaur's Roche Limit. Satellites could also form directly from the disrupted icy mantle. This formation mechanism predicts that roughly 10% of centaurs will have experienced potentially ring-forming encounters with giant planets. Ring systems of planets The composition of planetary ring particles varies, ranging from silicates to icy dust. Larger rocks and boulders may also be present, and in 2007 tidal effects from eight moonlets only a few hundred meters across were detected within Saturn's rings. The maximum size of a ring particle is determined by the specific strength of the material it is made of, its density, and the tidal force at its altitude. The tidal force is proportional to the average density inside the radius of the ring, or to the mass of the planet divided by the radius of the ring cubed. It is also inversely proportional to the square of the orbital period of the ring. Some planetary rings are influenced by shepherd moons, small moons that orbit near the inner or outer edges of a ringlet or within gaps in the rings. The gravity of shepherd moons serves to maintain a sharply defined edge to the ring; material that drifts closer to the shepherd moon's orbit is either deflected back into the body of the ring, ejected from the system, or accreted onto the moon itself. It is also predicted that Phobos, a moon of Mars, will break up and form into a planetary ring in about 50 million years. Its low orbit, with an orbital period that is shorter than a Martian day, is decaying due to tidal deceleration. Jupiter Jupiter's ring system was the third to be discovered, when it was first observed by the Voyager 1 probe in 1979, and was observed more thoroughly by the Galileo orbiter in the 1990s. Its four main parts are a faint thick torus known as the "halo"; a thin, relatively bright main ring; and two wide, faint "gossamer rings". The system consists mostly of dust. Saturn Saturn's rings are the most extensive ring system of any planet in the Solar System, and thus have been known to exist for quite some time. Galileo Galilei first observed them in 1610, but they were not accurately described as a disk around Saturn until Christiaan Huygens did so in 1655. The rings are not a series of tiny ringlets as many think, but are more of a disk with varying density. They consist mostly of water ice and trace amounts of rock, and the particles range in size from micrometers to meters. Uranus Uranus's ring system lies between the level of complexity of Saturn's vast system and the simpler systems around Jupiter and Neptune. They were discovered in 1977 by James L. Elliot, Edward W. Dunham, and Jessica Mink. In the time between then and 2005, observations by Voyager 2 and the Hubble Space Telescope led to a total of 13 distinct rings being identified, most of which are opaque and only a few kilometers wide. They are dark and likely consist of water ice and some radiation-processed organics. The relative lack of dust is due to aerodynamic drag from the extended exosphere-corona of Uranus. Neptune The system around Neptune consists of five principal rings that, at their densest, are comparable to the low-density regions of Saturn's rings. However, they are faint and dusty, much more similar in structure to those of Jupiter. The very dark material that makes up the rings is likely organics processed by radiation, like in the rings of Uranus. 20 to 70 percent of the rings are dust, a relatively high proportion. Hints of the rings were seen for decades prior to their conclusive discovery by Voyager 2 in 1989. Prehistoric ring systems Earth A 2024 study suggests that Earth may have had a ring system for a period of 40 million years, starting from the middle of the Ordovician period (around 466 million years ago). This ring system may have originated from a large asteroid that passed by Earth at this time and had a significant amount of debris stripped by Earth's gravitational pull, forming a ring system. Evidence for this ring comes from impact craters from the Ordovician meteor event appearing to cluster in a distinctive band around the Earth's equator at that time. The presence of this ring may have led to significant shielding of Earth from sun's rays and a severe cooling event, thus causing the Hirnantian glaciation, the coldest known period of the last 450 million years. Rings systems of minor planets and moons Reports in March 2008 suggested that Saturn's moon Rhea may have its own tenuous ring system, which would make it the only moon known to have a ring system. A later study published in 2010 revealed that imaging of Rhea by the Cassini spacecraft was inconsistent with the predicted properties of the rings, suggesting that some other mechanism is responsible for the magnetic effects that had led to the ring hypothesis. Prior to the arrival of New Horizons, some astronomers hypothesized that Pluto and Charon might have a circumbinary ring system created from dust ejected off of Pluto's small outer moons in impacts. A dust ring would have posed a considerable risk to the New Horizons spacecraft. However, this possibility was ruled out when New Horizons failed to detect any dust rings around Pluto. Chariklo 10199 Chariklo, a centaur, was the first minor planet discovered to have rings. It has two rings, perhaps due to a collision that caused a chain of debris to orbit it. The rings were discovered when astronomers observed Chariklo passing in front of the star UCAC4 248-108672 on June 3, 2013 from seven locations in South America. While watching, they saw two dips in the star's apparent brightness just before and after the occultation. Because this event was observed at multiple locations, the conclusion that the dip in brightness was in fact due to rings is unanimously the leading hypothesis. The observations revealed what is likely a -wide ring system that is about 1,000 times closer than the Moon is to Earth. In addition, astronomers suspect there could be a moon orbiting amidst the ring debris. If these rings are the leftovers of a collision as astronomers suspect, this would give fodder to the idea that moons (such as the Moon) form through collisions of smaller bits of material. Chariklo's rings have not been officially named, but the discoverers have nicknamed them Oiapoque and Chuí, after two rivers near the northern and southern ends of Brazil. Chiron A second centaur, 2060 Chiron, has a constantly evolving disk of rings. Based on stellar-occultation data that were initially interpreted as resulting from jets associated with Chiron's comet-like activity, the rings are proposed to be in radius, though their evolution does change the radius somewhat. Their changing appearance at different viewing angles can explain the long-term variation in Chiron's brightness over time. Chiron's rings are suspected to be maintained by orbiting material ejected during seasonal outbursts, as a third partial ring detected in 2018 had become a full ring by 2022, with an outburst in between in 2021. Haumea A ring around Haumea, a dwarf planet and resonant Kuiper belt member, was revealed by a stellar occultation observed on 21 January 2017. This makes it the first trans-Neptunian object found to have a ring system. The ring has a radius of about , a width of ≈ and an opacity of 0.5. The ring plane coincides with Haumea's equator and the orbit of its larger, outer moon Hi’iaka (which has a semimajor axis of ≈). The ring is close to the 3:1 resonance with Haumea's rotation, which is located at a radius of . It is well within Haumea's Roche limit, which would lie at a radius of about if Haumea were spherical (being nonspherical pushes the limit out farther). Quaoar In 2023, astronomers announced the discovery of a widely separated ring around the dwarf planet and Kuiper belt object Quaoar. Further analysis of the occultation data uncovered a second inner, fainter ring. Both rings display unusual properties. The outer ring orbits at a distance of , approximately 7.5 times the radius of Quaoar and more than double the distance of its Roche limit. The inner ring orbits at a distance of , approximately 4.6 times the radius of Quaoar and also beyond its Roche limit. The outer ring appears to be inhomogeneous, containing a thin, dense section as well as a broader, more diffuse section. Rings around exoplanets Because all giant planets of the Solar System have rings, the existence of exoplanets with rings is plausible. Although particles of ice, the material that is predominant in the rings of Saturn, can only exist around planets beyond the frost line, within this line rings consisting of rocky material can be stable in the long term. Such ring systems can be detected for planets observed by the transit method by additional reduction of the light of the central star if their opacity is sufficient. As of 2024, two candidate extrasolar ring systems have been found by this method, around HIP 41378 f and K2-33b. Fomalhaut b was found to be large and unclearly defined when detected in 2008. This was hypothesized to either be due to a cloud of dust attracted from the dust disc of the star, or a possible ring system, though in 2020 Fomalhaut b itself was determined to very likely be an expanding debris cloud from a collision of asteroids rather than a planet. Similarly, Proxima Centauri c has been observed to be far brighter than expected for its low mass of 7 Earth masses, which may be attributed to a ring system of about 5 . A 56-day-long sequence of dimming events in the star V1400 Centauri observed in 2007 was interpreted as a substellar object with a circumstellar disk or massive rings transiting the star. This substellar object, dubbed "J1407b", is most likely a free-floating brown dwarf or rogue planet several times the mass of Jupiter. The circumstellar disk or ring system of J1407b is about in radius. J1407b's transit of V1400 Centauri revealed gaps and density variations within its disk or ring system, which has been interpreted as hints of exomoons or exoplanets forming around J1407b. Visual comparison	Physical sciences	Planetary science	null
24731	https://en.wikipedia.org/wiki/Positron	Positron	The positron or antielectron is the particle with an electric charge of +1e, a spin of 1/2 (the same as the electron), and the same mass as an electron. It is the antiparticle (antimatter counterpart) of the electron. When a positron collides with an electron, annihilation occurs. If this collision occurs at low energies, it results in the production of two or more photons. Positrons can be created by positron emission radioactive decay (through weak interactions), or by pair production from a sufficiently energetic photon which is interacting with an atom in a material. History Theory In 1928, Paul Dirac published a paper proposing that electrons can have both a positive and negative charge. This paper introduced the Dirac equation, a unification of quantum mechanics, special relativity, and the then-new concept of electron spin to explain the Zeeman effect. The paper did not explicitly predict a new particle but did allow for electrons having either positive or negative energy as solutions. Hermann Weyl then published a paper discussing the mathematical implications of the negative energy solution. The positive-energy solution explained experimental results, but Dirac was puzzled by the equally valid negative-energy solution that the mathematical model allowed. Quantum mechanics did not allow the negative energy solution to simply be ignored, as classical mechanics often did in such equations; the dual solution implied the possibility of an electron spontaneously jumping between positive and negative energy states. However, no such transition had yet been observed experimentally. Dirac wrote a follow-up paper in December 1929 that attempted to explain the unavoidable negative-energy solution for the relativistic electron. He argued that "... an electron with negative energy moves in an external [electromagnetic] field as though it carries a positive charge." He further asserted that all of space could be regarded as a "sea" of negative energy states that were filled, so as to prevent electrons jumping between positive energy states (negative electric charge) and negative energy states (positive charge). The paper also explored the possibility of the proton being an island in this sea, and that it might actually be a negative-energy electron. Dirac acknowledged that the proton having a much greater mass than the electron was a problem, but expressed "hope" that a future theory would resolve the issue. Robert Oppenheimer argued strongly against the proton being the negative-energy electron solution to Dirac's equation. He asserted that if it were, the hydrogen atom would rapidly self-destruct. Weyl in 1931 showed that the negative-energy electron must have the same mass as that of the positive-energy electron. Persuaded by Oppenheimer's and Weyl's argument, Dirac published a paper in 1931 that predicted the existence of an as-yet-unobserved particle that he called an "anti-electron" that would have the same mass and the opposite charge as an electron and that would mutually annihilate upon contact with an electron. Richard Feynman, and earlier Ernst Stueckelberg, proposed an interpretation of the positron as an electron moving backward in time, reinterpreting the negative-energy solutions of the Dirac equation. Electrons moving backward in time would have a positive electric charge. John Archibald Wheeler invoked this concept to explain the identical properties shared by all electrons, suggesting that "they are all the same electron" with a complex, self-intersecting worldline. Yoichiro Nambu later applied it to all production and annihilation of particle-antiparticle pairs, stating that "the eventual creation and annihilation of pairs that may occur now and then is no creation or annihilation, but only a change of direction of moving particles, from the past to the future, or from the future to the past." The backwards in time point of view is nowadays accepted as completely equivalent to other pictures, but it does not have anything to do with the macroscopic terms "cause" and "effect", which do not appear in a microscopic physical description. Experimental clues and discovery Several sources have claimed that Dmitri Skobeltsyn first observed the positron long before 1930, or even as early as 1923. They state that while using a Wilson cloud chamber in order to study the Compton effect, Skobeltsyn detected particles that acted like electrons but curved in the opposite direction in an applied magnetic field, and that he presented photographs with this phenomenon in a conference in the University of Cambridge, on 23–27 July 1928. In his book on the history of the positron discovery from 1963, Norwood Russell Hanson has given a detailed account of the reasons for this assertion, and this may have been the origin of the myth. But he also presented Skobeltsyn's objection to it in an appendix. Later, Skobeltsyn rejected this claim even more strongly, calling it "nothing but sheer nonsense". Skobeltsyn did pave the way for the eventual discovery of the positron by two important contributions: adding a magnetic field to his cloud chamber (in 1925), and by discovering charged particle cosmic rays, for which he is credited in Carl David Anderson's Nobel lecture. Skobeltzyn did observe likely positron tracks on images taken in 1931, but did not identify them as such at the time. Likewise, in 1929 Chung-Yao Chao, a Chinese graduate student at Caltech, noticed some anomalous results that indicated particles behaving like electrons, but with a positive charge, though the results were inconclusive and the phenomenon was not pursued. Fifty years later, Anderson acknowledged that his discovery was inspired by the work of his Caltech classmate Chung-Yao Chao, whose research formed the foundation from which much of Anderson's work developed but was not credited at the time. Anderson discovered the positron on 2 August 1932, for which he won the Nobel Prize for Physics in 1936. Anderson did not coin the term positron, but allowed it at the suggestion of the Physical Review journal editor to whom he submitted his discovery paper in late 1932. The positron was the first evidence of antimatter and was discovered when Anderson allowed cosmic rays to pass through a cloud chamber and a lead plate. A magnet surrounded this apparatus, causing particles to bend in different directions based on their electric charge. The ion trail left by each positron appeared on the photographic plate with a curvature matching the mass-to-charge ratio of an electron, but in a direction that showed its charge was positive. Anderson wrote in retrospect that the positron could have been discovered earlier based on Chung-Yao Chao's work, if only it had been followed up on. Frédéric and Irène Joliot-Curie in Paris had evidence of positrons in old photographs when Anderson's results came out, but they had dismissed them as protons. The positron had also been contemporaneously discovered by Patrick Blackett and Giuseppe Occhialini at the Cavendish Laboratory in 1932. Blackett and Occhialini had delayed publication to obtain more solid evidence, so Anderson was able to publish the discovery first. Natural production Positrons are produced, together with neutrinos naturally in β+ decays of naturally occurring radioactive isotopes (for example, potassium-40) and in interactions of gamma quanta (emitted by radioactive nuclei) with matter. Antineutrinos are another kind of antiparticle produced by natural radioactivity (β− decay). Many different kinds of antiparticles are also produced by (and contained in) cosmic rays. In research published in 2011 by the American Astronomical Society, positrons were discovered originating above thunderstorm clouds; positrons are produced in gamma-ray flashes created by electrons accelerated by strong electric fields in the clouds. Antiprotons have also been found to exist in the Van Allen Belts around the Earth by the PAMELA module. Antiparticles, of which the most common are antineutrinos and positrons due to their low mass, are also produced in any environment with a sufficiently high temperature (mean particle energy greater than the pair production threshold). During the period of baryogenesis, when the universe was extremely hot and dense, matter and antimatter were continually produced and annihilated. The presence of remaining matter, and absence of detectable remaining antimatter, also called baryon asymmetry, is attributed to CP-violation: a violation of the CP-symmetry relating matter to antimatter. The exact mechanism of this violation during baryogenesis remains a mystery. Positron production from radioactive decay can be considered both artificial and natural production, as the generation of the radioisotope can be natural or artificial. Perhaps the best known naturally-occurring radioisotope which produces positrons is potassium-40, a long-lived isotope of potassium which occurs as a primordial isotope of potassium. Even though it is a small percentage of potassium (0.0117%), it is the single most abundant radioisotope in the human body. In a human body of mass, about 4,400 nuclei of 40K decay per second. The activity of natural potassium is 31 Bq/g. About 0.001% of these 40K decays produce about 4000 natural positrons per day in the human body. These positrons soon find an electron, undergo annihilation, and produce pairs of 511 keV photons, in a process similar (but much lower intensity) to that which happens during a PET scan nuclear medicine procedure. Recent observations indicate black holes and neutron stars produce vast amounts of positron-electron plasma in astrophysical jets. Large clouds of positron-electron plasma have also been associated with neutron stars. Observation in cosmic rays Satellite experiments have found evidence of positrons (as well as a few antiprotons) in primary cosmic rays, amounting to less than 1% of the particles in primary cosmic rays. However, the fraction of positrons in cosmic rays has been measured more recently with improved accuracy, especially at much higher energy levels, and the fraction of positrons has been seen to be greater in these higher energy cosmic rays. These do not appear to be the products of large amounts of antimatter from the Big Bang, or indeed complex antimatter in the universe (evidence for which is lacking, see below). Rather, the antimatter in cosmic rays appear to consist of only these two elementary particles. Recent theories suggest the source of such positrons may come from annihilation of dark matter particles, acceleration of positrons to high energies in astrophysical objects, and production of high energy positrons in the interactions of cosmic ray nuclei with interstellar gas. Preliminary results from the presently operating Alpha Magnetic Spectrometer (AMS-02) on board the International Space Station show that positrons in the cosmic rays arrive with no directionality, and with energies that range from 0.5 GeV to 500 GeV. Positron fraction peaks at a maximum of about 16% of total electron+positron events, around an energy of 275 ± 32 GeV. At higher energies, up to 500 GeV, the ratio of positrons to electrons begins to fall again. The absolute flux of positrons also begins to fall before 500 GeV, but peaks at energies far higher than electron energies, which peak about 10 GeV. These results on interpretation have been suggested to be due to positron production in annihilation events of massive dark matter particles. Positrons, like anti-protons, do not appear to originate from any hypothetical "antimatter" regions of the universe. On the contrary, there is no evidence of complex antimatter atomic nuclei, such as antihelium nuclei (i.e., anti-alpha particles), in cosmic rays. These are actively being searched for. A prototype of the AMS-02 designated AMS-01, was flown into space aboard the on STS-91 in June 1998. By not detecting any antihelium at all, the AMS-01 established an upper limit of 1.1×10−6 for the antihelium to helium flux ratio. Artificial production Physicists at the Lawrence Livermore National Laboratory in California have used a short, ultra-intense laser to irradiate a millimeter-thick gold target and produce more than 100 billion positrons. Presently significant lab production of 5 MeV positron-electron beams allows investigation of multiple characteristics such as how different elements react to 5 MeV positron interactions or impacts, how energy is transferred to particles, and the shock effect of gamma-ray bursts. In 2023, a collaboration between CERN and University of Oxford performed an experiment at the HiRadMat facility in which nano-second duration beams of electron-positron pairs were produced containing more than 10 trillion electron-positron pairs, so creating the first 'pair plasma' in the laboratory with sufficient density to support collective plasma behavior. Future experiments offer the possibility to study physics relevant to extreme astrophysical environments where copious electron-positron pairs are generated, such as gamma-ray bursts, fast radio bursts and blazar jets. Applications Certain kinds of particle accelerator experiments involve colliding positrons and electrons at relativistic speeds. The high impact energy and the mutual annihilation of these matter/antimatter opposites create a fountain of diverse subatomic particles. Physicists study the results of these collisions to test theoretical predictions and to search for new kinds of particles. The ALPHA experiment combines positrons with antiprotons to study properties of antihydrogen. Gamma rays, emitted indirectly by a positron-emitting radionuclide (tracer), are detected in positron emission tomography (PET) scanners used in hospitals. PET scanners create detailed three-dimensional images of metabolic activity within the human body. An experimental tool called positron annihilation spectroscopy (PAS) is used in materials research to detect variations in density, defects, displacements, or even voids, within a solid material.	Physical sciences	Antimatter	null
24733	https://en.wikipedia.org/wiki/Phencyclidine	Phencyclidine	Phencyclidine or phenylcyclohexyl piperidine (PCP), also known in its use as a street drug as angel dust among other names, is a dissociative anesthetic mainly used recreationally for its significant mind-altering effects. PCP may cause hallucinations, distorted perceptions of sounds, and violent behavior. As a recreational drug, it is typically smoked, but may be taken by mouth, snorted, or injected. It may also be mixed with cannabis or tobacco. Adverse effects may include seizures, coma, addiction, and an increased risk of suicide. Flashbacks may occur despite stopping usage. Chemically, PCP is a member of the arylcyclohexylamine class, and pharmacologically, it is a dissociative anesthetic. PCP works primarily as an NMDA receptor antagonist. PCP is most commonly used in the US. While usage peaked in the US in the 1970s, between 2005 and 2011, an increase in visits to emergency departments as a result of the drug occurred. As of 2022, in the US, about 0.7% of 12th grade students reported using PCP in the prior year, while 1.7% of people in the US over age 25 reported using it at some point in their lives. Recreational uses Phencyclidine is used for its ability to induce a dissociative state. Effects Behavioral effects can vary by dosage. Low doses produce numbness in the extremities and intoxication, characterized by staggering, unsteady gait, slurred speech, bloodshot eyes, and loss of balance. Moderate doses (5–10 mg intranasal, or 0.01–0.02 mg/kg intramuscular or intravenous) will produce analgesia and anesthesia. High doses may lead to convulsions. The drug is often illegally produced under poorly controlled conditions; this means that users may be unaware of the actual dose they are taking. Psychological effects include severe changes in body image, loss of ego boundaries, paranoia, and depersonalization. Psychosis, agitation and dysphoria, hallucinations, blurred vision, euphoria, and suicidal impulses are also reported, as well as occasional aggressive behavior. Like many other drugs, PCP has been known to alter mood states unpredictably, causing some individuals to become detached, and others to become animated. PCP may induce feelings of strength, power, and invulnerability as well as a numbing effect on the mind. Studies by the Drug Abuse Warning Network in the 1970s show that media reports of PCP-induced violence are greatly exaggerated and that incidents of violence are unusual and often limited to individuals with reputations for aggression regardless of drug use. Although uncommon, events of PCP-intoxicated individuals acting in an unpredictable fashion, possibly driven by their delusions or hallucinations, have been publicized. Other commonly cited types of incidents include inflicting property damage and self-mutilation of various types, such as pulling out one's teeth. These effects were not noted in its medicinal use in the 1950s and 1960s, however, reports of physical violence on PCP have often been shown to be unfounded. Recreational doses of the drug also occasionally appear to induce a psychotic state, with emotional and cognitive impairment that resembles a schizophrenic episode. Users generally report feeling detached from reality. Symptoms are summarized by the mnemonic device RED DANES: rage, erythema (redness of skin), dilated pupils, delusions, amnesia, nystagmus (oscillation of the eyeball when moving laterally), excitation, and skin dryness. Addiction PCP is self-administered and induces ΔFosB expression in the D1-type medium spiny neurons of the nucleus accumbens, and accordingly, excessive PCP use is known to cause addiction. PCP's rewarding and reinforcing effects are at least partly mediated by blocking the NMDA receptors in the glutamatergic inputs to D1-type medium spiny neurons in the nucleus accumbens. PCP has been shown to produce conditioned place aversion and conditioned place preference in animal studies. Schizophrenia A 2019 review found that the transition rate from a diagnosis of hallucinogen-induced psychosis (which included PCP) to that of schizophrenia was 26%. This was lower than cannabis-induced psychosis (34%) but higher than amphetamine (22%), opioid (12%), alcohol (10%), and sedative (9%) induced psychoses. In comparison, the transition rate to schizophrenia for "brief, atypical and not otherwise specified" psychosis was found to be 36%. Methods of administration PCP has multiple routes of administration. Most commonly, the powder form of the drug is snorted. PCP can also be orally ingested, injected subcutaneously or intravenously, or smoked laced with marijuana or cigarettes. PCP can be ingested through smoking. "Fry" and "sherm" are street terms for marijuana or tobacco cigarettes that are dipped in PCP and then dried. PCP hydrochloride can be insufflated (snorted), depending upon the purity. This is most often referred to as "angel dust". An oral pill can also be compressed from the co-compounded powder form of the drug. This is usually referred to as "peace pill". The free base is hydrophobic and may be absorbed through skin and mucous membranes (often inadvertently). This form of the drug is commonly called "wack". Management of intoxication Management of PCP intoxication mostly consists of supportive care – controlling breathing, circulation, and body temperature – and, in the early stages, treating psychiatric symptoms. Benzodiazepines, such as lorazepam, are the drugs of choice to control agitation and seizures (when present). Typical antipsychotics such as phenothiazines and haloperidol have been used to control psychotic symptoms, but may produce many undesirable side effects – such as dystonia – and their use is therefore no longer preferred; phenothiazines are particularly risky, as they may lower the seizure threshold, worsen hyperthermia, and boost the anticholinergic effects of PCP. If an antipsychotic is given, intramuscular haloperidol has been recommended. Forced acid diuresis (with ammonium chloride or, more safely, ascorbic acid) may increase the clearance of PCP from the body, and was somewhat controversially recommended in the past as a decontamination measure. However, it is now known that only around 10% of a dose of PCP is removed by the kidneys, which would make increased urinary clearance of little consequence; furthermore, urinary acidification is dangerous, as it may induce acidosis and worsen rhabdomyolysis (muscle breakdown), a not-unusual manifestation of PCP toxicity. Pharmacology Pharmacodynamics PCP is well known for its primary action on the NMDA receptor, an ionotropic glutamate receptor. As such, PCP is a non-competitive NMDA receptor antagonist. The role of NMDAR antagonism in the effect of PCP, ketamine, and related dissociative agents was first published in the early 1980s by David Lodge and colleagues. Other NMDA receptor antagonists include ketamine, tiletamine, dextromethorphan, nitrous oxide, and dizocilpine (MK-801). Research also indicates that PCP inhibits nicotinic acetylcholine receptors (nAChRs) among other mechanisms. Analogues of PCP exhibit varying potency at nACh receptors and NMDA receptors. Findings demonstrate that presynaptic nAChRs and NMDA receptor interactions influence the postsynaptic maturation of glutamatergic synapses and consequently impact synaptic development and plasticity in the brain. These effects can lead to inhibition of excitatory glutamate activity in certain brain regions such as the hippocampus and cerebellum thus potentially leading to memory loss as one of the effects of prolonged use. Acute effects on the cerebellum manifest as changes in blood pressure, breathing rate, pulse rate, and loss of muscular coordination during intoxication. PCP, like ketamine, also acts as a potent dopamine D2High receptor partial agonist in rat brain homogenate and has affinity for the human cloned D2High receptor. This activity may be associated with some of the other more psychotic features of PCP intoxication, which is evidenced by the successful use of D2 receptor antagonists (such as haloperidol) in the treatment of PCP psychosis. In addition to its well explored interactions with NMDA receptors, PCP has also been shown to inhibit dopamine reuptake, and thereby leads to increased extracellular levels of dopamine and hence increased dopaminergic neurotransmission. However, PCP has little affinity for the human monoamine transporters, including the dopamine transporter (DAT). Instead, its inhibition of monoamine reuptake may be mediated by interactions with allosteric sites on the monoamine transporters. PCP is notably a high-affinity ligand of the PCP site 2 (Ki = 154 nM), a not-well-characterized site associated with monoamine reuptake inhibition. Studies on rats indicate that PCP interacts indirectly with opioid receptors (endorphin and enkephalin) to produce analgesia. A binding study assessed PCP at 56 sites including neurotransmitter receptors and transporters and found that PCP had Ki values of >10,000 nM at all sites except the dizocilpine (MK-801) site of the NMDA receptor (Ki = 59 nM), the σ2 receptor (PC12) (Ki = 136 nM), and the serotonin transporter (Ki = 2,234 nM). The study notably found Ki values of >10,000 nM for the D2 receptor, the opioid receptors, the σ1 receptor, and the dopamine and norepinephrine transporters. These results suggest that PCP is a highly selective ligand of the NMDAR and σ2 receptor. However, PCP may also interact with allosteric sites on the monoamine transporters to produce inhibition of monoamine reuptake. Mechanism of action Phencyclidine is a noncompetitive NMDA receptor antagonist that blocks the activity of the NMDA receptor to cause anaesthesia and analgesia without causing cardiorespiratory depression. NMDA is an excitatory receptor in the brain, when activated normally the receptor acts as an ion channel and there is an influx of positive ions through the channel to cause nerve cell depolarisation. Phencyclidine inhibits the NMDA receptor by binding to the specific PCP binding site located within the ion channel. The PCP binding site is within close proximity to the magnesium blocking site, which may explain the similar inhibitory effects. Binding at the PCP site is mediated by two non-covalent interactions within the receptor: hydrogen bonding and hydrophobic interaction. Binding is also controlled by the gating mechanism of the ion channel. Because the PCP site is located within the ion channel, a coagonist such as glycine must bind and open the channel for PCP to enter, bind to the PCP site, and block the channel. Neurotoxicity Some studies found that, like other NMDA receptor antagonists, PCP can cause a kind of brain damage called Olney's lesions in rats. Studies conducted on rats showed that high doses of the NMDA receptor antagonist dizocilpine caused reversible vacuoles to form in certain regions of the rats' brains. All studies of Olney's lesions have only been performed on non-human animals and may not apply to humans. One unpublished study by Frank Sharp reportedly showed no damage by the NMDA antagonist ketamine, a structurally similar drug, far beyond recreational doses, but due to the study never having been published, its validity is controversial. PCP has also been shown to cause schizophrenia-like changes in N-acetylaspartate and N-acetylaspartylglutamate levels in the rat brain, which are detectable both in living rats and upon necropsy examination of brain tissue. It also induces symptoms in humans that mimic schizophrenia. PCP not only produced symptoms similar to schizophrenia, it also yielded electroencephalogram changes in the thalamocortical pathway (increased delta decreased alpha) and in the hippocampus (increase theta bursts) that were similar to those in schizophrenia. PCP-induced augmentation of dopamine release may link the NMDA and dopamine hypotheses of schizophrenia. Pharmacokinetics PCP is both water and lipid-soluble and is therefore distributed throughout the body quickly. PCP is metabolized into PCHP, PPC and PCAA. The drug is metabolized 90% by oxidative hydroxylation in the liver during the first pass. Metabolites are glucuronidated and excreted in the urine. Nine percent of ingested PCP is excreted in its unchanged form. When smoked, some of the compound is broken down by heat into 1-phenylcyclohexene (PC) and piperidine. The time taken before the effects of PCP manifest is dependent on the route of administration. The onset of action for inhalation occurs in 2–5 minutes, whereas the effects may take 15 to 60 minutes when ingested orally. Chemistry PCP is an arylcyclohexylamine. Analogues Fewer than 30 different analogs of PCP were reported as being used as a street drug during the 1970s and 1980s, mainly in the United States. Only of a few of these compounds were widely used including rolicyclidine (PCPy), eticyclidine (PCE), and tenocyclidine (TCP). Less common analogs include 3-HO-PCP, 3-MeO-PCMo, and 3-MeO-PCP. The generalized structural motif required for PCP-like activity is derived from structure-activity relationship studies of PCP derivatives. All of these derivatives are likely to share some of their psychoactive effects with PCP itself, although a range of potencies and varying mixtures of anesthetic, dissociative, and stimulant effects are known, depending on the particular drug and its substituents. In the United States, all of these compounds would be considered controlled substance analogs of PCP under the Federal Analog Act and are hence illegal drugs if sold for human consumption. History Phencyclidine was initially discovered in 1926 by and his student Paul Merkel as a product of a Grignard reaction of 1-piperidinocyclohexancarbonitrile. It was again synthesized in 1956 by chemist H Victor Maddox and brought to market as an anesthetic medication by pharmaceutical company Parke-Davis, now a subsidiary of Pfizer. Its use in humans was disallowed in the US in 1965 due to the high rates of side effects, while its use in animals was disallowed in 1978. Moreover, ketamine was discovered and was better tolerated as an anesthetic. PCP is classified as a schedule II drug in the US. Derivatives of PCP have been sold for recreational and non-medical use. Society and culture Regulation PCP is a Schedule II substance in the US. The Administrative Controlled Substances Code Number (ACSCN) for PCP is 7471. Its manufacturing quota for 2014 was . It is a Schedule I drug by the Controlled Drugs and Substances act in Canada, a List I drug of the Opium Law in the Netherlands, and a Class A substance in the UK. Frequency of use PCP began to emerge as a recreational drug in major cities in the US in the 1960s. In 1978, People magazine and Mike Wallace of the TV news program 60 Minutes called PCP the country's "number one drug problem". Although recreational use of the drug had always been relatively low, it began declining significantly in the 1980s. In surveys, the number of high school students admitting to trying PCP at least once fell from 13% in 1979 to less than 3% in 1990. Cultural depictions Jean-Michel Basquiat depicted two angel dust users in his 1982 painting Dustheads. Tsukasa Hojo's 1985 manga City Hunter features a drug, Angel Dust, presumably a reference to PCP's street name. The related 2023 animated film, City Hunter: Angel Dust, more directly moved the franchise's angel dust into the realm of fantasy, as it is portrayed as a science fiction nanomachine serum developed by a biotech company to create super-soldiers with a tendency to drive them berserk, side-stepping the real-life PCP. In Vivienne Medrano's adult animated musical comedy television series Hazbin Hotel, Angel Dust is an adult film star in Hell and one of the main protagonists, who in Hell took on the name "Angel Dust" as his chosen all-encompassing persona name, and one he uses exclusively in place of his actual name. It is intended as multipurpose for both his drag queen persona and his sex-work persona.	Biology and health sciences	Recreational drugs	Health
24749	https://en.wikipedia.org/wiki/Permian%E2%80%93Triassic%20extinction%20event	Permian–Triassic extinction event	Approximately 251.9 million years ago, the Permian–Triassic (P–T, P–Tr) extinction event (PTME; also known as the Late Permian extinction event, the Latest Permian extinction event, the End-Permian extinction event, and colloquially as the Great Dying) forms the boundary between the Permian and Triassic geologic periods, and with them the Paleozoic and Mesozoic eras. It is Earth's most severe known extinction event, with the extinction of 57% of biological families, 83% of genera, 81% of marine species and 70% of terrestrial vertebrate species. It is also the greatest known mass extinction of insects. It is the greatest of the "Big Five" mass extinctions of the Phanerozoic. There is evidence for one to three distinct pulses, or phases, of extinction. The scientific consensus is that the main cause of the extinction was the flood basalt volcanic eruptions that created the Siberian Traps, which released sulfur dioxide and carbon dioxide, resulting in euxinia (oxygen-starved, sulfurous oceans), elevating global temperatures, and acidifying the oceans. The level of atmospheric carbon dioxide rose from around 400 ppm to 2,500 ppm with approximately 3,900 to 12,000 gigatonnes of carbon being added to the ocean-atmosphere system during this period. Several other contributing factors have been proposed, including the emission of carbon dioxide from the burning of oil and coal deposits ignited by the eruptions; emissions of methane from the gasification of methane clathrates; emissions of methane by novel methanogenic microorganisms nourished by minerals dispersed in the eruptions; longer and more intense El Niño events; and an extraterrestrial impact which created the Araguainha crater and caused seismic release of methane and the destruction of the ozone layer with increased exposure to solar radiation. Dating Previously, it was thought that rock sequences spanning the Permian–Triassic boundary were too few and contained too many gaps for scientists to reliably determine its details. However, it is now possible to date the extinction with millennial precision. U–Pb zircon dates from five volcanic ash beds from the Global Stratotype Section and Point for the Permian–Triassic boundary at Meishan, China, establish a high-resolution age model for the extinction – allowing exploration of the links between global environmental perturbation, carbon cycle disruption, mass extinction, and recovery at millennial timescales. The first appearance of the conodont Hindeodus parvus has been used to delineate the Permian-Triassic boundary. The extinction occurred between and years ago, a duration of 60 ± 48 thousand years. A large, abrupt global decrease in δ13C, the ratio of the stable isotope carbon-13 to that of carbon-12, coincides with this extinction, and is sometimes used to identify the Permian–Triassic boundary and PTME in rocks that are unsuitable for radiometric dating. The negative carbon isotope excursion's magnitude was 4-7% and lasted for approximately 500 kyr, though estimating its exact value is challenging due to diagenetic alteration of many sedimentary facies spanning the boundary. Further evidence for environmental change around the Permian-Triassic boundary suggests an rise in temperature, and an increase in levels to (for comparison, the concentration immediately before the Industrial Revolution was , and the amount today is about 422 ppm). There is also evidence of increased ultraviolet radiation reaching the earth, causing the mutation of plant spores. It has been suggested that the Permian–Triassic boundary is associated with a sharp increase in the abundance of marine and terrestrial fungi, caused by the sharp increase in the amount of dead plants and animals fed upon by the fungi. This "fungal spike" has been used by some paleontologists to identify a lithological sequence as being on or very close to the Permian–Triassic boundary in rocks that are unsuitable for radiometric dating or have a lack of suitable index fossils. However, even the proposers of the fungal spike hypothesis pointed out that "fungal spikes" may have been a repeating phenomenon created by the post-extinction ecosystem during the earliest Triassic. The very idea of a fungal spike has been criticized on several grounds, including: Reduviasporonites, the most common supposed fungal spore, may be a fossilized alga; the spike did not appear worldwide; and in many places it did not fall on the Permian–Triassic boundary. The Reduviasporonites may even represent a transition to a lake-dominated Triassic world rather than an earliest Triassic zone of death and decay in some terrestrial fossil beds. Newer chemical evidence agrees better with a fungal origin for Reduviasporonites, diluting these critiques. Uncertainty exists regarding the duration of the overall extinction and about the timing and duration of various groups' extinctions within the greater process. Some evidence suggests that there were multiple extinction pulses or that the extinction was long and spread out over a few million years, with a sharp peak in the last million years of the Permian. Statistical analyses of some highly fossiliferous strata in Meishan, Zhejiang Province in southeastern China, suggest that the main extinction was clustered around one peak, while a study of the Liangfengya section found evidence of two extinction waves, MEH-1 and MEH-2, which varied in their causes, and a study of the Shangsi section showed two extinction pulses with different causes too. Recent research shows that different groups became extinct at different times; for example, while difficult to date absolutely, ostracod and brachiopod extinctions were separated by around 670,000 to 1.17 million years. Palaeoenvironmental analysis of Lopingian strata in the Bowen Basin of Queensland indicates numerous intermittent periods of marine environmental stress from the middle to late Lopingian leading up to the end-Permian extinction proper, supporting aspects of the gradualist hypothesis. Additionally, the decline in marine species richness and the structural collapse of marine ecosystems may have been decoupled as well, with the former preceding the latter by about 61,000 years according to one study. Whether the terrestrial and marine extinctions were synchronous or asynchronous is another point of controversy. Evidence from a well-preserved sequence in east Greenland suggests that the terrestrial and marine extinctions began simultaneously. In this sequence, the decline of animal life is concentrated in a period approximately 10,000 to 60,000 years long, with plants taking an additional several hundred thousand years to show the full impact of the event. Many sedimentary sequences from South China show synchronous terrestrial and marine extinctions. Research in the Sydney Basin of the PTME's duration and course also supports a synchronous occurrence of the terrestrial and marine biotic collapses. Other scientists believe the terrestrial mass extinction began between 60,000 and 370,000 years before the onset of the marine mass extinction. Chemostratigraphic analysis from sections in Finnmark and Trøndelag shows the terrestrial floral turnover occurred before the large negative δ13C shift during the marine extinction. Dating of the boundary between the Dicynodon and Lystrosaurus assemblage zones in the Karoo Basin indicates that the terrestrial extinction occurred earlier than the marine extinction. The Sunjiagou Formation of South China also records a terrestrial ecosystem demise predating the marine crisis. Other research still has found that the terrestrial extinction occurred after the marine extinction in the tropics. Studies of the timing and causes of the Permian-Triassic extinction are complicated by the often-overlooked Capitanian extinction (also called the Guadalupian extinction), just one of perhaps two mass extinctions in the late Permian that closely preceded the Permian-Triassic event. In short, when the Permian-Triassic starts it is difficult to know whether the end-Capitanian had finished, depending on the factor considered. Many of the extinctions once dated to the Permian-Triassic boundary have more recently been redated to the end-Capitanian. Further, it is unclear whether some species who survived the prior extinction(s) had recovered well enough for their final demise in the Permian-Triassic event to be considered separate from Capitanian event. A minority point of view considers the sequence of environmental disasters to have effectively constituted a single, prolonged extinction event, perhaps depending on which species is considered. This older theory, still supported in some recent papers, proposes that there were two major extinction pulses 9.4 million years apart, separated by a period of extinctions that were less extensive, but still well above the background level, and that the final extinction killed off only about 80% of marine species alive at that time, whereas the other losses occurred during the first pulse or the interval between pulses. According to this theory, one of these extinction pulses occurred at the end of the Guadalupian epoch of the Permian. For example, all dinocephalian genera died out at the end of the Guadalupian, as did the Verbeekinidae, a family of large-size fusuline foraminifera. The impact of the end-Guadalupian extinction on marine organisms appears to have varied between locations and between taxonomic groups – brachiopods and corals had severe losses. Extinction patterns Marine organisms Marine invertebrates suffered the greatest losses during the P–Tr extinction. Evidence of this was found in samples from south China sections at the P–Tr boundary. Here, 286 out of 329 marine invertebrate genera disappear within the final two sedimentary zones containing conodonts from the Permian. The decrease in diversity was probably caused by a sharp increase in extinctions, rather than a decrease in speciation. The extinction primarily affected organisms with calcium carbonate skeletons, especially those reliant on stable CO2 levels to produce their skeletons. These organisms were susceptible to the effects of the ocean acidification that resulted from increased atmospheric CO2. Organisms that relied on haemocyanin or haemoglobin for transporting oxygen were more resistant to extinction than those utilising haemerythrin or oxygen diffusion. There is also evidence that endemism was a strong risk factor influencing a taxon's likelihood of extinction. Bivalve taxa that were endemic and localised to a specific region were more likely to go extinct than cosmopolitan taxa. There was little latitudinal difference in the survival rates of taxa. Organisms that inhabited refugia less affected by global warming experienced lesser or delayed extinctions. Among benthic organisms the extinction event multiplied background extinction rates, and therefore caused maximum species loss to taxa that had a high background extinction rate (by implication, taxa with a high turnover). The extinction rate of marine organisms was catastrophic. Bioturbators were extremely severely affected, as evidenced by the loss of the sedimentary mixed layer in many marine facies during the end-Permian extinction. Surviving marine invertebrate groups included articulate brachiopods (those with a hinge), which had undergone a slow decline in numbers since the P–Tr extinction; the Ceratitida order of ammonites; and crinoids ("sea lilies"), which very nearly became extinct but later became abundant and diverse. The groups with the highest survival rates generally had active control of circulation, elaborate gas exchange mechanisms, and light calcification; more heavily calcified organisms with simpler breathing apparatuses suffered the greatest loss of species diversity. In the case of the brachiopods, at least, surviving taxa were generally small, rare members of a formerly diverse community. Conodonts were severely affected both in terms of taxonomic and morphological diversity, although not as severely as during the Capitanian mass extinction. The ammonoids, which had been in a long-term decline for the 30 million years since the Roadian (middle Permian), suffered a selective extinction pulse 10 million years before the main event, at the end of the Capitanian stage. In this preliminary extinction, which greatly reduced disparity, or the range of different ecological guilds, environmental factors were apparently responsible. Diversity and disparity fell further until the P–Tr boundary; the extinction here (P–Tr) was non-selective, consistent with a catastrophic initiator. During the Triassic, diversity rose rapidly, but disparity remained low. The range of morphospace occupied by the ammonoids, that is, their range of possible forms, shapes or structures, became more restricted as the Permian progressed. A few million years into the Triassic, the original range of ammonoid structures was once again reoccupied, but the parameters were now shared differently among clades. Ostracods experienced prolonged diversity perturbations during the Changhsingian before the PTME proper, when immense proportions of them abruptly vanished. At least 74% of ostracods died out during the PTME itself. Bryozoans had been on a long-term decline throughout the Late Permian epoch before they suffered even more catastrophic losses during the PTME, being the most severely affected clade among the lophophorates. Deep water sponges suffered a significant diversity loss and exhibited a decrease in spicule size over the course of the PTME. Shallow water sponges were affected much less strongly; they experienced an increase in spicule size and much lower loss of morphological diversity compared to their deep water counterparts. Foraminifera suffered a severe bottleneck in diversity. Evidence from South China indicates the foraminiferal extinction had two pulses. Foraminiferal biodiversity hotspots shifted into deeper waters during the PTME. Approximately 93% of latest Permian foraminifera became extinct, with 50% of the clade Textulariina, 92% of Lagenida, 96% of Fusulinida, and 100% of Miliolida disappearing. Foraminifera that were calcaerous suffered an extinction rate of 91%. The reason why lagenides survived while fusulinoidean fusulinides went completely extinct may have been due to the greater range of environmental tolerance and greater geographic distribution of the former compared to the latter. Cladodontomorph sharks likely survived the extinction by surviving in refugia in the deep oceans, a hypothesis based on the discovery of Early Cretaceous cladodontomorphs in deep, outer shelf environments. Ichthyosaurs, which evolved immediately before the PTME, were also PTME survivors. The Lilliput effect, the phenomenon of dwarfing of species during and immediately following a mass extinction event, has been observed across the Permian-Triassic boundary, notably occurring in foraminifera, brachiopods, bivalves, and ostracods. Though gastropods that survived the cataclysm were smaller in size than those that did not, it remains debated whether the Lilliput effect truly took hold among gastropods. Some gastropod taxa, termed "Gulliver gastropods", ballooned in size during and immediately following the mass extinction, exemplifying the Lilliput effect's opposite, which has been dubbed the Brobdingnag effect. Terrestrial invertebrates The Permian had great diversity in insect and other invertebrate species, including the largest insects ever to have existed. The end-Permian is the largest known mass extinction of insects; according to some sources, it may well be the only mass extinction to significantly affect insect diversity. Eight or nine insect orders became extinct and ten more were greatly reduced in diversity. Palaeodictyopteroids (insects with piercing and sucking mouthparts) began to decline during the mid-Permian; these extinctions have been linked to a change in flora. The greatest decline occurred in the Late Permian and was probably not directly caused by weather-related floral transitions. However, some observed entomofaunal declines in the PTME were biogeographic changes rather than outright extinctions. Terrestrial plants The geological record of terrestrial plants is sparse and based mostly on pollen and spore studies. Floral changes across the Permian-Triassic boundary are highly variable depending on the location and preservation quality of any given site. Plants are relatively immune to mass extinction, with the impact of all the major mass extinctions "insignificant" at a family level. Floral diversity losses were more superficial than those of marine animals. Even the reduction observed in species diversity (of 50%) may be mostly due to taphonomic processes. However, a massive rearrangement of ecosystems does occur, with plant abundances and distributions changing profoundly and all the forests virtually disappearing. The dominant floral groups changed, with many groups of land plants entering abrupt decline, such as Cordaites (gymnosperms) and Glossopteris (seed ferns). The severity of plant extinction has been disputed. The Glossopteris-dominated flora that characterised high-latitude Gondwana collapsed in Australia around 370,000 years before the Permian-Triassic boundary, with this flora's collapse being less constrained in western Gondwana but still likely occurring a few hundred thousand years before the boundary. The collapse of this flora is indirectly marked by an abrupt change in river morphology from meandering to braided river systems, signifying the widespread demise of rooted plants. Palynological or pollen studies from East Greenland of sedimentary rock strata laid down during the extinction period indicate dense gymnosperm woodlands before the event. At the same time that marine invertebrate macrofauna declined, these large woodlands died out and were followed by a rise in diversity of smaller herbaceous plants including Lycopodiophyta, both Selaginellales and Isoetales. Data from Kap Stosch suggest that floral species richness was not significantly affected during the PTME. The Cordaites flora, which dominated the Angaran floristic realm corresponding to Siberia, collapsed over the course of the extinction. In the Kuznetsk Basin, the aridity-induced extinction of the regions's humid-adapted forest flora dominated by cordaitaleans occurred approximately 252.76 Ma, around 820,000 years before the end-Permian extinction in South China, suggesting that the end-Permian biotic catastrophe may have started earlier on land and that the ecological crisis may have been more gradual and asynchronous on land compared to its more abrupt onset in the marine realm. In North China, the transition between the Upper Shihhotse and Sunjiagou Formations and their lateral equivalents marked a very large extinction of plants in the region. Those plant genera that did not go extinct still experienced a great reduction in their geographic range. Following this transition, coal swamps vanished. The North Chinese floral extinction correlates with the decline of the Gigantopteris flora of South China. In South China, the subtropical Cathaysian gigantopterid dominated rainforests abruptly collapsed. The floral extinction in South China is associated with bacterial blooms in soil and nearby lacustrine ecosystems, with soil erosion resulting from the die-off of plants being their likely cause. Wildfires too likely played a role in the fall of Gigantopteris. A conifer flora in what is now Jordan, known from fossils near the Dead Sea, showed unusual stability over the Permian-Triassic transition, and appears to have been only minimally affected by the crisis. Terrestrial vertebrates The tempo of the terrestrial vertebrate extinction is disputed. Some evidence from the Karoo Basin indicates a protracted extinction lasting a million years. Other evidence from the Karoo deposits suggest it took 50,000 years or less, while a study of coprolites in the Vyazniki fossil beds in Russia suggests it took only a few thousand years. Aridification induced by global warming was the chief culprit behind terrestrial vertebrate extinctions. There is enough evidence to indicate that over two thirds of terrestrial labyrinthodont amphibians, sauropsid ("reptile") and therapsid ("proto-mammal") taxa became extinct. Large herbivores suffered the heaviest losses. All Permian anapsid reptiles died out except the procolophonids (although testudines have morphologically-anapsid skulls, they are now thought to have separately evolved from diapsid ancestors). Pelycosaurs died out before the end of the Permian. Too few Permian diapsid fossils have been found to support any conclusion about the effect of the Permian extinction on diapsids (the "reptile" group from which lizards, snakes, crocodilians, and dinosaurs (including birds) evolved). Tangasaurids were largely unaffected. Gorgonopsians are traditionally thought to have gone extinct during the PTME, but some tentative evidence suggests they may have survived into the Triassic. Freshwater and euryhaline fishes, having experienced minimal diversity losses before the PTME, were unaffected during the PTME and actually appear to have increased in diversity across the Permian-Triassic boundary. However, faunal turnovers in freshwater fish communities occurred in areas like the Kuznetsk Basin. The groups that survived suffered extremely heavy losses of species and some terrestrial vertebrate groups very nearly became extinct at the end of the Permian. Some of the surviving groups did not persist for long past this period, but others that barely survived went on to produce diverse and long-lasting lineages. However, it took 30million years for the terrestrial vertebrate fauna to fully recover both numerically and ecologically. It is difficult to analyze extinction and survival rates of land organisms in detail because few terrestrial fossil beds span the Permian–Triassic boundary. The best-known record of vertebrate changes across the Permian–Triassic boundary occurs in the Karoo Supergroup of South Africa, but statistical analyses have so far not produced clear conclusions. One study of the Karoo Basin found that 69% of terrestrial vertebrates went extinct over 300,000 years leading up to the Permian-Triassic boundary, followed by a minor extinction pulse involving four taxa that survived the previous extinction interval. Another study of latest Permian vertebrates in the Karoo Basin found that 54% of them went extinct due to the PTME. Biotic recovery In the wake of the extinction event, the ecological structure of present-day biosphere evolved from the stock of surviving taxa. In the sea, the "Palaeozoic evolutionary fauna" declined while the "modern evolutionary fauna" achieved greater dominance; the Permian-Triassic mass extinction marked a key turning point in this ecological shift that began after the Capitanian mass extinction and culminated in the Late Jurassic. Typical taxa of shelly benthic faunas were now bivalves, snails, sea urchins and Malacostraca, whereas bony fishes and marine reptiles diversified in the pelagic zone. On land, dinosaurs and mammals arose in the course of the Triassic. The profound change in the taxonomic composition was partly a result of the selectivity of the extinction event, which affected some taxa (e.g., brachiopods) more severely than others (e.g., bivalves). However, recovery was also differential between taxa. Some survivors became extinct some million years after the extinction event without having rediversified (dead clade walking, e.g. the snail family Bellerophontidae), whereas others rose to dominance over geologic times (e.g., bivalves). Marine ecosystems A cosmopolitanism event began immediately after the end-Permian extinction event. Marine post-extinction faunas were mostly species-poor and were dominated by few disaster taxa such as the bivalves Claraia, Unionites, Eumorphotis, and Promyalina, the conodonts Clarkina and Hindeodus, the inarticulate brachiopod Lingularia, and the foraminifera Earlandia and Rectocornuspira kalhori, the latter of which is sometimes classified under the genus Ammodiscus. Their guild diversity was also low. Post-PTME faunas had a flat, insignificant latitudinal diversity gradient. The speed of recovery from the extinction is disputed. Some scientists estimate that it took 10 million years (until the Middle Triassic) due to the severity of the extinction. However, studies in Bear Lake County, near Paris, Idaho, and nearby sites in Idaho and Nevada showed a relatively quick rebound in a localized Early Triassic marine ecosystem (Paris biota), taking around 1.3 million years to recover, while an unusually diverse and complex ichnobiota is known from Italy less than a million years after the end-Permian extinction. Additionally, the complex Guiyang biota found near Guiyang, China also indicates life thrived in some places just a million years after the mass extinction, as does a fossil assemblage known as the Shanggan fauna found in Shanggan, China, the Wangmo biota from the Luolou Formation of Guizhou, and a gastropod fauna from the Al Jil Formation of Oman. Regional differences in the pace of biotic recovery existed, which suggests that the impact of the extinction may have been felt less severely in some areas than others, with differential environmental stress and instability being the source of the variance. In addition, it has been proposed that although overall taxonomic diversity rebounded rapidly, functional ecological diversity took much longer to return to its pre-extinction levels; one study concluded that marine ecological recovery was still ongoing 50 million years after the extinction, during the latest Triassic, even though taxonomic diversity had rebounded in a tenth of that time. The pace and timing of recovery also differed based on clade and mode of life. Seafloor communities maintained a comparatively low diversity until the end of the Early Triassic, approximately 4 million years after the extinction event. Epifaunal benthos took longer to recover than infaunal benthos. This slow recovery stands in remarkable contrast with the quick recovery seen in nektonic organisms such as ammonoids, which exceeded pre-extinction diversities already two million years after the crisis, and conodonts, which diversified considerably over the first two million years of the Early Triassic. Recent work suggests that the pace of recovery was intrinsically driven by the intensity of competition among species, which drives rates of niche differentiation and speciation. That recovery was slow in the Early Triassic can be explained by low levels of biological competition due to the paucity of taxonomic diversity, and that biotic recovery explosively accelerated in the Anisian can be explained by niche crowding, a phenomenon that would have drastically increased competition, becoming prevalent by the Anisian. Biodiversity rise thus behaved as a positive feedback loop enhancing itself as it took off in the Spathian and Anisian. Accordingly, low levels of interspecific competition in seafloor communities that are dominated by primary consumers correspond to slow rates of diversification and high levels of interspecific competition among nektonic secondary and tertiary consumers to high diversification rates. Other explanations state that life was delayed in its recovery because grim conditions returned periodically over the course of the Early Triassic, causing further extinction events, such as the Smithian-Spathian boundary extinction. Continual episodes of extremely hot climatic conditions during the Early Triassic have been held responsible for the delayed recovery of oceanic life, in particular skeletonised taxa that are most vulnerable to high carbon dioxide concentrations. The relative delay in the recovery of benthic organisms has been attributed to widespread anoxia, but high abundances of benthic species contradict this explanation. A 2019 study attributed the dissimilarity of recovery times between different ecological communities to differences in local environmental stress during the biotic recovery interval, with regions experiencing persistent environmental stress post-extinction recovering more slowly, supporting the view that recurrent environmental calamities were culpable for retarded biotic recovery. Recurrent Early Triassic environmental stresses also acted as a ceiling limiting the maximum ecological complexity of marine ecosystems until the Spathian. Recovery biotas appear to have been ecologically uneven and unstable into the Anisian, making them vulnerable to environmental stresses. Whereas most marine communities were fully recovered by the Middle Triassic, global marine diversity reached pre-extinction values no earlier than the Middle Jurassic, approximately 75 million years after the extinction event. Prior to the extinction, about two-thirds of marine animals were sessile and attached to the seafloor. During the Mesozoic, only about half of the marine animals were sessile while the rest were free-living. Analysis of marine fossils from the period indicated a decrease in the abundance of sessile epifaunal suspension feeders such as brachiopods and sea lilies and an increase in more complex mobile species such as snails, sea urchins and crabs. Before the Permian mass extinction event, both complex and simple marine ecosystems were equally common. After the recovery from the mass extinction, the complex communities outnumbered the simple communities by nearly three to one, and the increase in predation pressure and durophagy led to the Mesozoic Marine Revolution. Marine vertebrates recovered relatively quickly, with complex predator-prey interactions with vertebrates at the top of the food web being known from coprolites five million years after the PTME. Post-PTME hybodonts exhibited extremely rapid tooth replacement. Ichthyopterygians appear to have ballooned in size extremely rapidly following the PTME. Bivalves rapidly recolonised many marine environments in the wake of the catastrophe. Bivalves were fairly rare before the P–Tr extinction but became numerous and diverse in the Triassic, taking over niches that were filled primarily by brachiopods before the mass extinction event. Bivalves were once thought to have outcompeted brachiopods, but this outdated hypothesis about the brachiopod-bivalve transition has been disproven by Bayesian analysis. The success of bivalves in the aftermath of the extinction event may have been a function of them possessing greater resilience to environmental stress compared to the brachiopods that they coexisted with, whilst other studies have emphasised the greater niche breadth of the former. The rise of bivalves to taxonomic and ecological dominance over brachiopods was not synchronous, however, and brachiopods retained an outsized ecological dominance into the Middle Triassic even as bivalves eclipsed them in taxonomic diversity. Some researchers think the brachiopod-bivalve transition was attributable not only to the end-Permian extinction but also the ecological restructuring that began as a result of the Capitanian extinction. Infaunal habits in bivalves became more common after the PTME. Linguliform brachiopods were commonplace immediately after the extinction event, their abundance having been essentially unaffected by the crisis. Adaptations for oxygen-poor and warm environments, such as increased lophophoral cavity surface, shell width/length ratio, and shell miniaturisation, are observed in post-extinction linguliforms. The surviving brachiopod fauna was very low in diversity and exhibited no provincialism whatsoever. Brachiopods began their recovery around 250.1 ± 0.3 Ma, as marked by the appearance of the genus Meishanorhynchia, believed to be the first of the progenitor brachiopods that evolved after the mass extinction. Major brachiopod rediversification only began in the late Spathian and Anisian in conjunction with the decline of widespread anoxia and extreme heat and the expansion of more habitable climatic zones. Brachiopod taxa during the Anisian recovery interval were only phylogenetically related to Late Permian brachiopods at a familial taxonomic level or higher; the ecology of brachiopods had radically changed from before in the mass extinction's aftermath. Ostracods were extremely rare during the basalmost Early Triassic. Taxa associated with microbialites were disproportionately represented among ostracod survivors. Ostracod recovery began in the Spathian. Despite high taxonomic turnover, the ecological life modes of Early Triassic ostracods remained rather similar to those of pre-PTME ostracods. Bryozoans in the Early Triassic were restricted to the Boreal realm. They were also not diverse, represented mainly by members of Trepostomatida. During the Middle Triassic, there was a rise in bryozoan diversity, which peaked in the Carnian. However, bryozoans took until the Late Cretaceous to recover their full diversity. Crinoids ("sea lilies") suffered a selective extinction, resulting in a decrease in the variety of their forms. Though cladistic analyses suggest the beginning of their recovery to have taken place in the Induan, the recovery of their diversity as measured by fossil evidence was far less brisk, showing up in the late Ladinian. Their adaptive radiation after the extinction event resulted in forms possessing flexible arms becoming widespread; motility, predominantly a response to predation pressure, also became far more prevalent. Though their taxonomic diversity remained relatively low, crinoids regained much of their ecological dominance by the Middle Triassic epoch. Stem-group echinoids survived the PTME. The survival of miocidarid echinoids such as Eotiaris is likely attributable to their ability to thrive in a wide range of environmental conditions. Conodonts saw a rapid recovery during the Induan, with anchignathodontids experiencing a diversity peak in the earliest Induan. Gondolellids diversified at the end of the Griesbachian; this diversity spike was most responsible for the overall conodont diversity peak in the Smithian. Segminiplanate conodonts again experienced a brief period of domination in the early Spathian, probably related to a transient oxygenation of deep waters. Neospathodid conodonts survived the crisis but underwent proteromorphosis. In the PTME's aftermath, disaster taxa of benthic foraminifera filled many of their vacant niches. The recovery of benthic foraminifera was very slow and frequently interrupted until the Spathian. In the Tethys, foraminiferal communities remained low in diversity into the Middle Triassic, with the exception of a notable Ladinian fauna from the Catalonian Basin. Microbial reefs were common across shallow seas for a short time during the earliest Triassic, predominating in low latitudes while being rarer in higher latitudes, occurring both in anoxic and oxic waters. Polybessurus-like microfossils often dominated these earliest Triassic microbialites. Microbial-metazoan reefs appeared very early in the Early Triassic; and they dominated many surviving communities across the recovery from the mass extinction. Microbialite deposits appear to have declined in the early Griesbachian synchronously with a significant sea level drop that occurred then. Metazoan-built reefs reemerged during the Olenekian, mainly being composed of sponge biostrome and bivalve builups. Keratose sponges were particularly noteworthy in their integral importance to Early Triassic microbial-metazoan reef communities, and they helped to create stability in heavily damaged ecosystems during early phases of biotic recovery. "Tubiphytes"-dominated reefs appeared at the end of the Olenekian, representing the earliest platform-margin reefs of the Triassic, though they did not become abundant until the late Anisian, when reefs' species richness increased. The first scleractinian corals appear in the late Anisian as well, although they would not become the dominant reef builders until the end of the Triassic period. Bryozoans, after sponges, were the most numerous organisms in Tethyan reefs during the Anisian. Metazoan reefs became common again during the Anisian because the oceans cooled down then from their overheated state during the Early Triassic. Biodiversity amongst metazoan reefs did not recover until well into the Anisian, millions of years after non-reef ecosystems recovered their diversity. Microbially induced sedimentary structures (MISS) from the earliest Triassic have been found to be associated with abundant opportunistic bivalves and vertical burrows, and it is likely that post-extinction microbial mats played a vital, indispensable role in the survival and recovery of various bioturbating organisms. The microbialite refuge hypothesis has been criticised as reflecting a taphonomic bias due to the greater preservation potential of microbialite deposits, however, rather than a genuine phenomenon. Ichnocoenoses show that marine ecosystems recovered to pre-extinction levels of ecological complexity by the late Olenekian. Anisian ichnocoenoses show slightly lower diversity than Spathian ichnocoenoses, although this was likely a taphonomic consequence of increased and deeper bioturbation erasing evidence of shallower bioturbation. Ichnological evidence suggests that recovery and recolonisation of marine environments may have taken place by way of outward dispersal from refugia that suffered relatively mild perturbations and whose local biotas were less strongly affected by the mass extinction compared to the rest of the world's oceans. Although complex bioturbation patterns were rare in the Early Triassic, likely reflecting the inhospitability of many shallow water environments in the extinction's wake, complex ecosystem engineering managed to persist locally in some places, and may have spread from there after harsh conditions across the global ocean were ameliorated over time. Wave-dominated shoreface settings (WDSS) are believed to have served as refugium environments because they appear to have been unusually diverse in the mass extinction's aftermath. Terrestrial plants The proto-recovery of terrestrial floras took place from a few tens of thousands of years after the end-Permian extinction to around 350,000 years after it, with the exact timeline varying by region. Furthermore, severe extinction pulses continued to occur after the Permian-Triassic boundary, causing additional floral turnovers. Gymnosperms recovered within a few thousand years after the Permian-Triassic boundary, but around 500,000 years after it, the Dominant gymnosperm genera were replaced by lycophytesextant lycophytes are recolonizers of disturbed areasduring an extinction pulse at the Griesbachian-Dienerian boundary. The particular post-extinction dominance of lycophytes, which were well adapted for coastal environments, can be explained in part by global marine transgressions during the Early Triassic. The worldwide recovery of gymnosperm forests took approximately 4–5 million years. However, this trend of prolonged lycophyte dominance during the Early Triassic was not universal, as evidenced by the much more rapid recovery of gymnosperms in certain regions, and floral recovery likely did not follow a congruent, globally universal trend but instead varied by region according to local environmental conditions. In East Greenland, lycophytes replaced gymnosperms as the dominant plants. Later, other groups of gymnosperms again become dominant but again suffered major die-offs. These cyclical flora shifts occurred a few times over the course of the extinction period and afterward. These fluctuations of the dominant flora between woody and herbaceous taxa indicate chronic environmental stress resulting in a loss of most large woodland plant species. The successions and extinctions of plant communities do not coincide with the shift in values but occurred many years after. In what is now the Barents Sea of the coast of Norway, the post-extinction flora is dominated by pteridophytes and lycopods, which were suited for primary succession and recolonisation of devastated areas, although gymnosperms made a rapid recovery, with the lycopod dominated flora not persisting across most of the Early Triassic as postulated in other regions. In Europe and North China, the interval of recovery was dominated by the lycopsid Pleuromeia, an opportunistic pioneer plant that filled ecological vacancies until other plants were able to expand out of refugia and recolonise the land. Conifers became common by the early Anisian, while pteridosperms and cycadophytes only fully recovered by the late Anisian. During the survival phase in the terrestrial extinction's immediate aftermath, from the latest Changhsingian to the Griesbachian, South China was dominated by opportunistic lycophytes. Low-lying herbaceous vegetation dominated by the isoetalean Tomiostrobus was ubiquitous following the collapse of the gigantopterid-dominated forests of before. In contrast to the highly biodiverse gigantopterid rainforests, the post-extinction landscape of South China was near-barren and had vastly lower diversity. Plant survivors of the PTME in South China experienced extremely high rates of mutagenesis induced by heavy metal poisoning. From the late Griesbachian to the Smithian, conifers and ferns began to rediversify. After the Smithian, the opportunistic lycophyte flora declined, as the newly radiating conifer and fern species permanently replaced them as the dominant components of South China's flora. In Tibet, the early Dienerian Endosporites papillatus–Pinuspollenites thoracatus assemblages closely resemble late Changhsingian Tibetan floras, suggesting that the widespread, dominant latest Permian flora resurged easily after the PTME. However, in the late Dienerian, a major shift towards assemblages dominated by cavate trilete spores took place, heralding widespread deforestation and a rapid change to hotter, more humid conditions. Quillworts and spike mosses dominated Tibetan flora for about a million years after this shift. In Pakistan, then the northern margin of Gondwana, the flora was rich in lycopods associated with conifers and pteridosperms. Floral turnovers continued to occur due to repeated perturbations arising from recurrent volcanic activity until terrestrial ecosystems stabilised around 2.1 Myr after the PTME. In southwestern Gondwana, the post-extinction flora was dominated by bennettitaleans and cycads, with members of Peltaspermales, Ginkgoales, and Umkomasiales being less common constituents of this flora. Around the Induan-Olenekian boundary, as palaeocommunities recovered, a new Dicroidium flora was established, in which Umkomasiales continued to be prominent and in which Equisetales and Cycadales were subordinate forms. The Dicroidium flora further diversified in the Anisian to its peak, wherein Umkomasiales and Ginkgoales constituted most of the tree canopy and Peltaspermales, Petriellales, Cycadales, Umkomasiales, Gnetales, Equisetales, and Dipteridaceae dominated the understory. Coal gap No coal deposits are known from the Early Triassic, and those in the Middle Triassic are thin and low-grade. This "coal gap" has been explained in many ways. It has been suggested that new, more aggressive fungi, insects, and vertebrates evolved and killed vast numbers of trees. These decomposers themselves suffered heavy losses of species during the extinction and are not considered a likely cause of the coal gap. It could simply be that all coal-forming plants were rendered extinct by the P–Tr extinction and that it took 10 million years for a new suite of plants to adapt to the moist, acid conditions of peat bogs. Abiotic factors (factors not caused by organisms), such as decreased rainfall or increased input of clastic sediments, may also be to blame. On the other hand, the lack of coal may simply reflect the scarcity of all known sediments from the Early Triassic. Coal-producing ecosystems, rather than disappearing, may have moved to areas where we have no sedimentary record for the Early Triassic. For example, in eastern Australia a cold climate had been the norm for a long period, with a peat mire ecosystem adapted to these conditions. Approximately 95% of these peat-producing plants went locally extinct at the P–Tr boundary; coal deposits in Australia and Antarctica disappear significantly before the P–Tr boundary. Terrestrial vertebrates Land vertebrates took an unusually long time to recover from the P–Tr extinction; palaeontologist Michael Benton estimated the recovery was not complete until after the extinction, i.e. not until the Late Triassic, when the first dinosaurs had risen from bipedal archosaurian ancestors and the first mammals from small cynodont ancestors. A tetrapod gap may have existed from the Induan until the early Spathian between ~30 °N and ~ 40 °S due to extreme heat making these low latitudes uninhabitable for these animals. During the hottest phases of this interval, the gap would have spanned an even greater latitudinal range. East-central Pangaea, with its relatively wet climate, served as a dispersal corridor for PTME survivors during their Early Triassic recolonisation of the supercontinent. In North China, tetrapod body and ichnofossils are extremely rare in Induan facies, but become more abundant in the Olenekian and Anisian, showing a biotic recovery of tetrapods synchronous with the decreasing aridity during the Olenekian and Anisian. In Russia, even after 15 Myr of recovery, during which ecosystems were rebuilt and remodelled, many terrestrial vertebrate guilds were absent, including small insectivores, small piscivores, large herbivores, and apex predators. Coprolitic evidence indicates that freshwater food webs had recovered by the early Ladinian, with a lacustrine coprolite assemblage from the Ordos Basin of China providing evidence of a trophically multileveled ecosystem containing at least six different trophic levels. The highest trophic levels were filled by vertebrate predators. Overall, terrestrial faunas after the extinction event tended to be more variable and heterogeneous across space than those of the Late Permian, which exhibited less provincialism, being much more geographically homogeneous. Synapsids Lystrosaurus, a pig-sized herbivorous dicynodont therapsid, constituted as much as 90% of some earliest Triassic land vertebrate fauna, although some recent evidence has called into question its status as a post-PTME disaster taxon. The dicynodont genus is often used as a biostratigraphic marker for the PTME. The evolutionary success of Lystrosaurus in the aftermath of the PTME is believed to be attributable to the dicynodont taxon's grouping behaviour and tolerance for extreme and highly variable climatic conditions. Other likely factors behind the success of Lystrosaurus included extremely fast growth rate exhibited by the dicynodont genus, along with its early onset of sexual maturity. Antarctica may have served as a refuge for dicynodonts during the PTME from which surviving dicynodonts spread out of in its aftermath. Ichnological evidence from the earliest Triassic of the Karoo Basin shows dicynodonts were abundant in the immediate aftermath of the biotic crisis. Smaller carnivorous cynodont therapsids also survived, a group that included the ancestors of mammals. As with dicynodonts, selective pressures favoured endothermic epicynodonts. Therocephalians likewise survived; burrowing may have been a key adaptation that helped them make it through the PTME. In the Karoo region of southern Africa, the therocephalians Tetracynodon, Moschorhinus and Ictidosuchoides survived, but do not appear to have been abundant in the Triassic. Early Triassic therocephalians were mostly survivors of the PTME rather than newly evolved taxa that originated during the evolutionary radiation in its aftermath. Both therocephalians and cynodonts, known collectively as eutheriodonts, decreased in body size from the Late Permian to the Early Triassic. This decrease in body size has been interpreted as an example of the Lilliput effect. Sauropsids Archosaurs (which included the ancestors of dinosaurs and crocodilians) were initially rarer than therapsids, but they began to displace therapsids in the mid-Triassic. Olenekian tooth fossil assemblages from the Karoo Basin indicate that archosauromorphs were already highly diverse by this point in time, though not very ecologically specialised. In the mid to late Triassic, the dinosaurs evolved from one group of archosaurs, and went on to dominate terrestrial ecosystems during the Jurassic and Cretaceous. This "Triassic Takeover" may have contributed to the evolution of mammals by forcing the surviving therapsids and their mammaliform successors to live as small, mainly nocturnal insectivores; nocturnal life probably forced at least the mammaliforms to develop fur, better hearing and higher metabolic rates, while losing part of the differential color-sensitive retinal receptors reptilians and birds preserved. Archosaurs also experienced an increase in metabolic rates over time during the Early Triassic. The archosaur dominance would end again due to the Cretaceous–Paleogene extinction event, after which both birds (only extant dinosaurs) and mammals (only extant synapsids) would diversify and share the world. Temnospondyls Temnospondyl amphibians made a quick recovery; the appearance in the fossil record of so many temnospondyl clades suggests they may have been ideally suited as pioneer species that recolonised decimated ecosystems. During the Induan, tupilakosaurids in particular thrived as disaster taxa, including Tupilakosaurus itself, though they gave way to other temnospondyls as ecosystems recovered. Temnospondyls were reduced in size during the Induan, but their body size rebounded to pre-PTME levels during the Olenekian. Mastodonsaurus and trematosaurians were the main aquatic and semiaquatic predators during most of the Triassic, some preying on tetrapods and others on fish. Terrestrial invertebrates Most fossil insect groups found after the Permian–Triassic boundary differ significantly from those before: Of Paleozoic insect groups, only the Glosselytrodea, Miomoptera, and Protorthoptera have been discovered in deposits from after the extinction. The caloneurodeans, paleodictyopteroids, protelytropterans, and protodonates became extinct by the end of the Permian. Though Triassic insects are very different from those of the Permian, a gap in the insect fossil record spans approximately 15 million years from the late Permian to early Triassic. In well-documented Late Triassic deposits, fossils overwhelmingly consist of modern fossil insect groups. Microbially induced sedimentary structures (MISS) dominated North Chinese terrestrial fossil assemblages in the Early Triassic. In Arctic Canada as well, MISS became a common occurrence following the Permian-Triassic extinction. The prevalence of MISS in many Early Triassic rocks shows that microbial mats were an important feature of post-extinction ecosystems that were denuded of bioturbators that would have otherwise prevented their widespread occurrence. The disappearance of MISS later in the Early Triassic likely indicated a greater recovery of terrestrial ecosystems and specifically a return of prevalent bioturbation. Hypotheses about cause Explaining an event from 250 million years ago is inherently difficult, with much of the evidence on land eroded or deeply buried, while the spreading seafloor is completely recycled over 200 million years, leaving no useful indications beneath the ocean. Yet, scientists have gathered significant evidence for causes, and several mechanisms have been proposed. The proposals include both catastrophic and gradual processes (similar to those theorized for the Cretaceous–Paleogene extinction event, but with much less current consensus). The catastrophic group includes one or more large bolide impact events, increased volcanism, and sudden release of methane from the seafloor, either due to dissociation of methane hydrate deposits or metabolism of organic carbon deposits by methanogenic microbes. The gradual group includes sea level change, increasing hypoxia, and increasing aridity. Any hypothesis about the cause must explain the selectivity of the event, which affected organisms with calcium carbonate skeletons most severely; the long period (4 to 6 million years) before recovery started, and the minimal extent of biological mineralization (despite inorganic carbonates being deposited) once the recovery began. Volcanism Siberian Traps The flood basalt eruptions that produced the large igneous province of the Siberian Traps were among the largest known volcanic events, extruding lava over , roughly the size of Saudi Arabia, producing a catastrophic impact. The date of the Siberian Traps eruptions matches well with the extinction event. A study of the Norilsk and Maymecha-Kotuy regions of the northern Siberian platform indicates that volcanic activity occurred during a few enormous pulses of magma, as opposed to more regular flows. The Siberian Traps caused one of the most rapid rises of atmospheric carbon dioxide levels in the geologic record, with the rate of carbon dioxide emissions estimated as five times faster than during the preceding catastrophic Capitanian mass extinction during the eruption of the Emeishan Traps. Overwhelming inorganic carbon sinks, carbon dioxide levels might have jumped from between 500 and 4,000 ppm prior to the extinction to around 8,000 ppm after, according to one estimate. Another study estimated pre-extinction carbon dioxide levels at 400 ppm, which then rose to 2,500 ppm, with 3,900 to 12,000 gigatonnes of carbon added to the ocean-atmosphere system. Extreme temperature rise would have followed, though some evidence suggests a lag of 12,000 to 128,000 years between the rise in volcanic carbon dioxide emissions and global warming. During the latest Permian before the extinction, global average surface temperatures were about 18.2 °C, which shot up to as much as 35 °C, this hyperthermal condition lasting as long as 500,000 years. Air temperatures at Gondwana's high southern latitudes experienced a warming of ~10–14 °C. According to oxygen isotope shifts from conodont apatite in South China, low latitude surface water temperatures surged about 8 °C. In present-day Iran, tropical sea surface temperatures were between 27 and 33 °C during the Changhsingian but jumped to over 35 °C during the PTME. The increased mean state temperatures also brought stronger El Nino events, heightening short-term climate variability. These extremely high atmospheric carbon dioxide concentrations persisted over a long period. The position and alignment of Pangaea at the time made the inorganic carbon cycle very inefficient at burying carbon. In a 2020 paper, scientists reconstructed the mechanisms that led to the extinction event in a biogeochemical model, showed the consequences of the greenhouse effect on the marine environment, and concluded that the mass extinction can be traced back to volcanic CO emissions. Evidence also points to volcanic combustion of underground fossil fuel deposits, based on paired coronene-mercury spikes coinciding with geographically widespread mercury anomalies and the rise in isotopically light carbon. Te/Th values increase twentyfold over the PTME, further indicating it was concomitant with extreme volcanism. A major volcanogenic influx of isotopically light zinc from the Siberian Traps has also been recorded, further confirming that volcanism was contemporary with the PTME. The Siberian Traps eruptions had unusual features that made them even more dangerous. The Siberian lithosphere is rich in halogens extremely destructive to the ozone layer, and evidence from subcontinental lithospheric xenoliths indicates that as much as 70% of their halogen content was released into the atmosphere. Around 18 teratonnes of hydrochloric acid were emitted, along with sulphur-rich volatiles that caused dust clouds and acid aerosols, which would have blocked out sunlight and disrupted photosynthesis on land and in the photic zone of the ocean, causing food chains to collapse. These volcanic outbursts of sulphur also induced brief but severe global cooling punctuating the broader trend of rapid global warming, with glacio-eustatic sea level fall. However, the briefness of these cold events makes them unlikely to have been a significant kill mechanism. The eruptions may also have caused acid rain as the aerosols washed out of the atmosphere. That may have killed land plants and mollusks and planktonic organisms with calcium carbonate shells. Pure flood basalts produce fluid, low-viscosity lava, and do not hurl debris into the atmosphere. It appears, however, that 20% of the output of the Siberian Traps eruptions was pyroclastic ash thrown high into the atmosphere, increasing the short-term cooling effect. When this had washed out of the atmosphere, the excess carbon dioxide would have remained and global warming would have proceeded unchecked. Burning of hydrocarbon deposits may have exacerbated the extinction. The Siberian Traps are underlain by thick sequences of Early-Mid Paleozoic aged carbonate and evaporite deposits, as well as Carboniferous-Permian aged coal bearing clastic rocks. When heated, such as by igneous intrusions, these rocks may emit large amounts of greenhouse and toxic gases. The unique setting of the Siberian Traps over these deposits is likely the reason for the severity of the extinction. The basalt lava erupted or intruded into carbonate rocks and sediments in the process of forming large coal beds, which would have emitted large amounts of carbon dioxide, leading to stronger global warming after the dust and aerosols settled. The change of the eruptions from flood basalt to sill dominated emplacement, liberating even more trapped hydrocarbon deposits, coincides with the main onset of the extinction and is linked to a major negative excursion. The intermediate temperature of the Siberian Traps magmas optimised the extremely voluminous release of CO2 by way of heating of evaporites and carbonates. Venting of coal-derived methane was accompanied by explosive combustion of coal and discharge of coal-fly ash. A 2011 study led by Stephen E. Grasby reported evidence that volcanism caused massive coal beds to ignite, possibly releasing more than 3 trillion tons of carbon. They found ash deposits in deep rock layers near what is now the Buchanan Lake Formation: "coal ash dispersed by the explosive Siberian Trap eruption would be expected to have an associated release of toxic elements in impacted water bodies where fly ash slurries developed. ... Mafic megascale eruptions are long-lived events that would allow significant build-up of global ash clouds." Grasby said, "In addition to these volcanoes causing fires through coal, the ash it spewed was highly toxic and was released in the land and water, potentially contributing to the worst extinction event in earth history." However, some researchers propose that these supposed fly ashes were actually the result of wildfires not related to massive coal combustion by intrusive magmatism. A 2013 study led by Q.Y. Yang reported that the total amounts of important volatiles emitted from the Siberian Traps consisted of 8.5 × Tg CO, 4.4 × Tg CO, 7.0 × Tg HS, and 6.8 × Tg SO. The sill-dominated emplacement of the Siberian Traps prolonged their warming effects; whereas extrusive volcanism generates an abundance of subaerial basalts that efficiently sequester carbon dioxide via the silicate weathering process, underground sills cannot sequester atmospheric carbon dioxide and mitigate global warming. Additionally, enhanced reverse weathering and depletion of siliceous carbon sinks enabled extreme warmth to persist for much longer than expected if the excess carbon dioxide was sequestered by silicate rock. The reduction in marine primary productivity diminished emissions of dimethyl sulphate and dimethylsulphoniopropionate, enhancing warming. Also, the decline in biological silicate deposition resulting from the mass extinction of siliceous organisms acted as a positive feedback loop wherein mass death of marine life exacerbated and prolonged extreme hothouse conditions by depleting yet another siliceous carbon sink. Mercury anomalies corresponding to the time of Siberian Traps activity have been found in many geographically disparate sites, indicating that these volcanic eruptions released significant quantities of toxic mercury into the atmosphere and ocean, causing even larger terrestrial and marine die-offs. A series of surges raised terrestrial and marine environmental mercury concentrations by orders of magnitude above normal background levels and caused mercury poisoning over periods of a thousand years each. Mutagenesis in surviving plants after the PTME coeval with mercury and copper loading confirms volcanically induced heavy metal toxicity. Increased bioproductivity may have sequestered mercury and party mitigated poisoning. Immense volumes of nickel aerosols and cobalt and arsenic emisions, were also released, further contributing to metal poisoning. The devastation wrought by the Siberian Traps did not end following the Permian-Triassic boundary. Carbon isotope fluctuations suggest that massive Siberian Traps activity recurred many times during the Early Triassic, a finding corroborated by mercury spikes, causing further extinction events during the epoch. Choiyoi Silicic Large Igneous Province A second flood basalt event that produced the Choiyoi Silicic Large Igneous Province in southwestern Gondwana between around 286 Ma and 247 Ma has also been suggested as a significant additional extinction mechanism. At 1,300,000 cubic kilometres in volume and 1,680,000 square kilometres in area, this event was 40% the size of the Siberian Traps. Specifically, this flood basalt has been implicated in the regional demise of the Gondwanan Glossopteris flora. Indochina-South China subduction-zone volcanic arc Mercury anomalies preceding the end-Permian extinction have been discovered in what was then the boundary between the South China craton and the Indochinese plate, a subduction zone with a volcanic arc. Hafnium isotopes from syndepositional magmatic zircons found in ash beds created by this volcanic pulse confirm its origin in subduction-zone volcanism rather than large igneous province activity. The enrichment of copper samples from these deposits in isotopically light copper provide additional confirmation. This volcanism has been speculated to have caused local biotic stress among radiolarians, sponges, and brachiopods over the 60,000 years preceding the end-Permian marine extinction, as well as an ammonoid crisis with decreased morphological complexity and size and increased rate of turnover beginning in the lower C. yini biozone, around 200,000 years before the extinction. Methane clathrate gasification Methane clathrates, also known as methane hydrates, consist of molecules of methane trapped in the crystal lattice of ice. This methane, produced by methanogen microbes, has a about 6% below normal ( −6.0%). At the right combination of pressure and temperature, clathrates form near the surface of permafrost and in large quantities on continental shelves and nearby seabed at water depths of at least , buried in sediments up to below the sea floor. Massive release of methane from these clathrates may have contributed to the PTME, as scientists have found worldwide evidence of a swift decrease of about 1% in thein carbonate rocks from the end-Permian. This is the first, largest, and fastest of a series of excursions (decreases and increases) in the ratio, until it abruptly stabilised in the middle Triassic, followed soon afterwards by the recovery of calcifying shelled sealife. The seabed probably contained methane hydrate deposits, and the lava caused the deposits to dissociate, releasing vast quantities of methane. A vast release of methane might cause significant global warming since methane is a very powerful greenhouse gas. Strong evidence suggests the global temperatures increased by about 6 °C (10.8 °F) near the equator and therefore by more at higher latitudes: a sharp decrease in oxygen isotope ratios (); the extinction of Glossopteris flora (Glossopteris and plants that grew in the same areas), which needed a cold climate, with its replacement by floras typical of lower paleolatitudes. It was also suggested that a large-scale release of methane and other greenhouse gases from the ocean into the atmosphere was connected to the anoxic events and euxinic (sulfidic) events at the time, with the exact mechanism compared to the 1986 Lake Nyos disaster. The clathrate hypothesis seemed the only proposed mechanism sufficient to cause a global 1% reduction in the . While a variety of factors may have contributed to the ratio drop, a 2002 review found most of them insufficient to account for the observed amount: Gases from volcanic eruptions have aabout 0.5 to 0.8% below standard ( −0.5 to −0.8%), but a 1995 assessment concluded that the observed 1.0% worldwide reduction would have required eruptions massively larger than any found. (However, this analysis addressed only CO2 produced by the magma itself, not from interactions with carbon bearing sediments, as described below.) A reduction in organic activity would extract C more slowly from the environment and leave more of it to be incorporated into sediments, thus reducing the Biochemical processes preferentially use the lighter isotopes since chemical reactions are ultimately driven by electromagnetic forces between atoms and lighter isotopes respond more quickly to these forces, but a study of a smaller drop of 0.3 to 0.4% in ( −3 to −4 ‰) at the Paleocene-Eocene Thermal Maximum (PETM) concluded that even transferring all the organic carbon (in organisms, soils, and dissolved in the ocean) into sediments would be insufficient: Even such a large burial of material rich in C would not have produced the 'smaller' drop in the of the rocks around the PETM. Buried sedimentary organic matter has a 2.0 to 2.5% below normal ( −2.0 to −2.5%). Theoretically, if the sea level fell sharply, shallow marine sediments would be exposed to oxidation. But 6,500–8,400 gigatonnes (1 gigatonne = kg) of organic carbon would have to be oxidized and returned to the ocean-atmosphere system within less than a few hundred thousand years to reduce the by 1.0%, which is not thought to be a realistic possibility. Moreover, sea levels were rising rather than falling at the time of the extinction. Rather than a sudden decline in sea level, intermittent periods of ocean-bottom hyperoxia and anoxia (high-oxygen and low- or zero-oxygen conditions) may have caused the fluctuations in the Early Triassic; and global anoxia may have been responsible for the end-Permian blip. The continents of the end-Permian and early Triassic were more clustered in the tropics than they are now, and large tropical rivers would have dumped sediment into smaller, partially enclosed ocean basins at low latitudes. Such conditions favor oxic and anoxic episodes; oxic/anoxic conditions would result in a rapid release/burial, respectively, of large amounts of organic carbon, which has a low because biochemical processes use the lighter isotopes more. That or another organic-based reason may have been responsible for both that and a late Proterozoic/Cambrian pattern of fluctuating However, the clathrate hypothesis has also been criticized. Carbon-cycle models which include consideration of roasting carbonate sediments by volcanism confirm that it would have had enough effect to produce the observed reduction. Also, the pattern of isotope shifts expected to result from a massive release of methane does not match the patterns seen throughout the Early Triassic. Not only would such a cause require the release of five times as much methane as postulated for the PETM, but would it also have to be reburied at an unrealistically high rate to account for the rapid increases in the (episodes of high positive ) throughout the early Triassic before it was released several times again. The latest research suggests that greenhouse gas release during the extinction event was dominated by volcanic carbon dioxide, and while methane release had to have contributed, isotopic signatures show that thermogenic methane released from the Siberian Traps had consistently played a larger role than methane from clathrates and any other biogenic sources such as wetlands during the event. Adding to the evidence against methane clathrate release as the central driver of warming, the main rapid warming event is also associated with marine transgression rather than regression; the former would not normally have initiated methane release, which would have instead required a decrease in pressure, something that would be generated by a retreat of shallow seas. The configuration of the world's landmasses into one supercontinent would also mean that the global gas hydrate reservoir was lower than today, further damaging the case for methane clathrate dissolution as a major cause of the carbon cycle disruption. Hypercapnia and acidification Marine organisms are more sensitive to changes in (carbon dioxide) levels than terrestrial organisms for a variety of reasons. is 28 times more soluble in water than is oxygen. Marine animals normally function with lower concentrations of in their bodies than land animals, as the removal of in air-breathing animals is impeded by the need for the gas to pass through the respiratory system's membranes (lungs' alveolus, tracheae, and the like), even when diffuses more easily than oxygen. In marine organisms, relatively modest but sustained increases in concentrations hamper the synthesis of proteins, reduce fertilization rates, and produce deformities in calcareous hard parts. Higher concentrations of also result in decreased activity levels in many active marine animals, hindering their ability to obtain food. An analysis of marine fossils from the Permian's final Changhsingian stage found that marine organisms with a low tolerance for hypercapnia (high concentration of carbon dioxide) had high extinction rates, and the most tolerant organisms had very slight losses. The most vulnerable marine organisms were those that produced calcareous hard parts (from calcium carbonate) and had low metabolic rates and weak respiratory systems, notably calcareous sponges, rugose and tabulate corals, calcite-depositing brachiopods, bryozoans, and echinoderms; about 81% of such genera became extinct. Close relatives without calcareous hard parts suffered only minor losses, such as sea anemones, from which modern corals evolved. Animals with high metabolic rates, well-developed respiratory systems, and non-calcareous hard parts had negligible losses except for conodonts, in which 33% of genera died out. This pattern is also consistent with what is known about the effects of hypoxia, a shortage but not total absence of oxygen. However, hypoxia cannot have been the only killing mechanism for marine organisms. Nearly all of the continental shelf waters would have had to become severely hypoxic to account for the magnitude of the extinction, but such a catastrophe would make it difficult to explain the very selective pattern of the extinction. Mathematical models of the Late Permian and Early Triassic atmospheres show a significant but protracted decline in atmospheric oxygen levels, with no acceleration near the P–Tr boundary. Minimum atmospheric oxygen levels in the Early Triassic are never less than present-day levels and so the decline in oxygen levels does not match the temporal pattern of the extinction. In addition, an increase in concentration is inevitably linked to ocean acidification, consistent with the preferential extinction of heavily calcified taxa and other signals in the rock record that suggest a more acidic ocean, such as a carbonate production crisis that occurred a few thousand years after volcanic greenhouse gas emissions began. The decrease in ocean pH is calculated to be up to 0.7 units. An extreme aragonite sea formed. Ocean acidification was most extreme at mid-latitudes, and the major marine transgression associated with the end-Permian extinction is believed to have devastated shallow shelf communities in conjunction with anoxia. Evidence from paralic facies spanning the Permian-Triassic boundary in western Guizhou and eastern Yunnan, however, shows a local marine transgression dominated by carbonate deposition, suggesting that ocean acidification did not occur across the entire globe and was likely limited to certain regions of the world's oceans. One study, published in Scientific Reports, concluded that widespread ocean acidification, if it did occur, was not intense enough to impede calcification and only occurred during the beginning of the extinction event. The relative success of many marine organisms that were very vulnerable to acidification has further been used to argue that acidification was not a major extinction contributor. The persistence of highly elevated carbon dioxide concentrations in the atmosphere and ocean during the Early Triassic would have impeded the recovery of biocalcifying organisms after the PTME. Acidity generated by increased carbon dioxide concentrations in soil and sulphur dioxide dissolution in rainwater was also a kill mechanism on land. The increasing acidification of rainwater caused increased soil erosion as a result of the increased acidity of forest soils, evidenced by the increased influx of terrestrially derived organic sediments found in marine sedimentary deposits during the end-Permian extinction. Further evidence of an increase in soil acidity comes from elevated Ba/Sr ratios in earliest Triassic soils. A positive feedback loop further enhancing and prolonging soil acidification may have resulted from the decline of infaunal invertebrates like tubificids and chironomids, which remove acid metabolites from the soil. The increased abundance of vermiculitic clays in Shansi, South China coinciding with the Permian-Triassic boundary strongly suggests a sharp drop in soil pH causally related to volcanogenic emissions of carbon dioxide and sulphur dioxide. Hopane anomalies have also been interpreted as evidence of acidic soils and peats. As with many other environmental stressors, acidity on land episodically persisted well into the Triassic, stunting the recovery of terrestrial ecosystems. Anoxia and euxinia Evidence for widespread ocean anoxia (severe deficiency of oxygen) and euxinia (presence of hydrogen sulfide) is found from the Late Permian to the Early Triassic. Throughout most of the Tethys and Panthalassic Oceans, evidence for anoxia appears at the extinction event, including small pyrite framboids, negative δ238U excursions, negative δ15N excursions, positive δ82/78Se isotope excursions, relatively positive δ13C ratios in polycyclic aromatic hydrocarbons, high Th/U ratios, positive Ce/Ce* anomalies, depletions of molybdenum, uranium, and vanadium from seawater, and fine laminations in sediments. However, evidence for anoxia precedes the extinction at some other sites, including Spiti, India, Shangsi, China, Meishan, China, Opal Creek, Alberta, and Kap Stosch, Greenland. Biogeochemical evidence also points to the presence of euxinia during the PTME. Biomarkers for green sulfur bacteria, such as isorenieratane, the diagenetic product of isorenieratene, are widely used as indicators of photic zone euxinia because green sulfur bacteria require both sunlight and hydrogen sulfide to survive. Their abundance in sediments from the P–T boundary indicates euxinic conditions were present even in the shallow waters of the photic zone. Negative mercury isotope excursions further bolster evidence for extensive euxinia during the PTME. The disproportionate extinction of high-latitude marine species provides further evidence for oxygen depletion as a killing mechanism; low-latitude species living in warmer, less oxygenated waters are naturally better adapted to lower levels of oxygen and are able to migrate to higher latitudes during periods of global warming, whereas high-latitude organisms are unable to escape from warming, hypoxic waters at the poles. Evidence of a lag between volcanic mercury inputs and biotic turnovers provides further support for anoxia and euxinia as the key killing mechanism, because extinctions would be expected to be synchronous with volcanic mercury discharge if volcanism and hypercapnia was the primary driver of extinction. The sequence of extinctions in some sections, with deep water organisms being affected first followed by shallow water and then by bottom water organisms, is believed to reflect the migration of oxygen minimum zones. Models of ocean chemistry suggest that anoxia and euxinia were closely associated with hypercapnia. This suggests that poisoning from hydrogen sulfide, anoxia, and hypercapnia acted together as a killing mechanism. Hypercapnia best explains the selectivity of the extinction, but anoxia and euxinia probably contributed to the high mortality of the event. The sequence of events leading to anoxic oceans may have been triggered by carbon dioxide emissions from the eruption of the Siberian Traps. In that scenario, warming from the enhanced greenhouse effect would reduce the solubility of oxygen in seawater, causing the concentration of oxygen to decline. Increased coastal evaporation would have caused the formation of warm saline bottom water (WSBW) depleted in oxygen and nutrients, which spread across the world through the deep oceans. The influx of WSBW caused thermal expansion of water that raised sea levels, bringing anoxic waters onto shallow shelfs and enhancing the formation of WSBW in a positive feedback loop. The flux of terrigeneous material into the oceans increased as a result of soil erosion, which would have facilitated increased eutrophication; marine regression likewise enhanced terrigeneous material inputs. Increased chemical weathering of the continents due to warming and the acceleration of the water cycle would increase the riverine flux of nutrients to the ocean. Additionally, the Siberian Traps directly fertilised the oceans with iron and phosphorus as well, triggering bioblooms and marine snowstorms. Increased phosphate levels would have supported greater primary productivity in the surface oceans. The increase in organic matter production would have caused more organic matter to sink into the deep ocean, where its respiration would further decrease oxygen concentrations. Once anoxia became established, it would have been sustained by a positive feedback loop because deep water anoxia tends to increase the recycling efficiency of phosphate, leading to even higher productivity. Along the Panthalassan coast of South China, oxygen decline was also driven by large-scale upwelling of deep water enriched in various nutrients, causing this region of the ocean to be hit by especially severe anoxia. Convective overturn helped facilitate the expansion of anoxia throughout the water column. A severe anoxic event at the end of the Permian would have allowed sulfate-reducing bacteria to thrive, causing the production of large amounts of hydrogen sulfide in the anoxic ocean, turning it euxinic. In some regions, anoxia briefly disappeared when transient cold snaps resulting from volcanic sulphur emissions occurred. The persistence of anoxia through the Early Triassic may explain the slow recovery of marine life and low levels of biodiversity after the extinction, particularly that of benthic organisms. Anoxia disappeared from shallow waters more rapidly than the deep ocean. Reexpansions of oxygen-minimum zones did not cease until the late Spathian, periodically setting back and restarting the biotic recovery process. The decline in continental weathering towards the end of the Spathian at last began ameliorating marine life from recurrent anoxia. In some regions of Panthalassa, pelagic zone anoxia continued to recur as late as the Anisian, probably due to increased productivity and a return of aeolian upwelling. Some sections show a rather quick return to oxic water column conditions, however, so for how long anoxia persisted remains debated. The volatility of the Early Triassic sulphur cycle suggests marine life continued to face returns of euxinia as well. Some scientists have challenged the anoxia hypothesis on the grounds that long-lasting anoxic conditions could not have been supported if Late Permian thermohaline ocean circulation conformed to the "thermal mode" characterised by cooling at high latitudes. Anoxia may have persisted under a "haline mode" in which circulation was driven by subtropical evaporation, although the "haline mode" is highly unstable and was unlikely to have represented Late Permian oceanic circulation. Oxygen depletion via extensive microbial blooms also played a role in the biological collapse of not just marine ecosystems but freshwater ones as well. Persistent lack of oxygen after the extinction event itself helped delay biotic recovery for much of the Early Triassic epoch. Aridification Increasing continental aridity, a trend well underway even before the PTME as a result of the coalescence of the supercontinent Pangaea, was drastically exacerbated by terminal Permian volcanism and global warming. The combination of global warming and drying generated an increased incidence of wildfires. Tropical coastal swamp floras such as those in South China may have been very detrimentally impacted by the increase in wildfires, though it is ultimately unclear if an increase in wildfires played a role in driving taxa to extinction. Aridification trends varied widely in their tempo and regional impact. Analysis of the fossil river deposits of the floodplains of the Karoo Basin indicate a shift from meandering to braided river patterns, indicating a very abrupt drying of the climate. The climate change may have taken as little as 100,000 years, prompting the extinction of the unique Glossopteris flora and its associated herbivores, followed by the carnivorous guild. A pattern of aridity-induced extinctions that progressively ascended up the food chain has been deduced from Karoo Basin biostratigraphy. Evidence for aridification in the Karoo across the Permian-Triassic boundary is not, however, universal, as some palaeosol evidence indicates a wettening of the local climate during the transition between the two geologic periods. Evidence from the Sydney Basin of eastern Australia, on the other hand, suggests that the expansion of semi-arid and arid climatic belts across Pangaea was not immediate but was instead a gradual, prolonged process. Apart from the disappearance of peatlands, there was little evidence of significant sedimentological changes in depositional style across the Permian-Triassic boundary. Instead, a modest shift to amplified seasonality and hotter summers is suggested by palaeoclimatological models based on weathering proxies from the region's Late Permian and Early Triassic deposits. In the Kuznetsk Basin of southwestern Siberia, an increase in aridity led to the demise of the humid-adapted Cordaites forests in the region a few hundred thousand years before the Permian-Triassic boundary. Drying of this basin has been attributed to a broader poleward shift of drier, more arid climates during the late Changhsingian before the more abrupt main phase of the extinction at the Permian-Triassic boundary that disproportionately affected tropical and subtropical species. The persistence of hyperaridity varied regionally as well. In the North China Basin, highly arid climatic conditions are recorded during the latest Permian, near the Permian-Triassic boundary, with a swing towards increased precipitation during the Early Triassic, the latter likely assisting biotic recovery following the mass extinction. Elsewhere, such as in the Karoo Basin, episodes of dry climate recurred regularly in the Early Triassic, with profound effects on terrestrial tetrapods. The types and diversity of ichnofossils in a locality has been used as an indicator measuring aridity. Nurra, an ichnofossil site on the island of Sardinia, shows evidence of major drought-related stress among crustaceans. Whereas the Permian subnetwork at Nurra displays extensive horizontal backfilled traces and high ichnodiversity, the Early Triassic subnetwork is characterised by an absence of backfilled traces, an ichnological sign of aridification. Ozone depletion A collapse of the atmospheric ozone shield has been invoked as an explanation for the mass extinction, particularly that of terrestrial plants. Ozone production may have been reduced by 60-70%, increasing the flux of ultraviolet radiation by 400% at equatorial latitudes and 5,000% at polar latitudes. The hypothesis has the advantage of explaining the mass extinction of plants, which would have added to the methane levels and should otherwise have thrived in an atmosphere with a high level of carbon dioxide. Fossil spores from the end-Permian further support the theory; many spores show deformities that could have been caused by ultraviolet radiation, which would have been more intense after hydrogen sulfide emissions weakened the ozone layer. Malformed plant spores from the time of the extinction event show an increase in ultraviolet B absorbing compounds, confirming that increased ultraviolet radiation played a role in the environmental catastrophe and excluding other possible causes of mutagenesis, such as heavy metal toxicity, in these mutated spores. Extremely positive Δ33S anomalies provide evidence of photolysis of volcanic SO2, indicating increased ultraviolet radiation flux. Sulphur isotope data from North China are inconsistent with a total collapse of the ozone layer, however, suggesting it may have not been as major a kill mechanism as others. Multiple mechanisms could have reduced the ozone shield and rendered it ineffective. Computer modelling shows high atmospheric methane levels are associated with ozone shield decline and may have contributed to its reduction during the PTME. Volcanic emissions of sulphate aerosols into the stratosphere would have dealt significant destruction to the ozone layer. As mentioned previously, the rocks in the region where the Siberian Traps were emplaced are extremely rich in halogens. The intrusion of Siberian Traps volcanism into deposits rich in organohalogens, such as methyl bromide and methyl chloride, would have been another source of ozone destruction. An uptick in wildfires, a natural source of methyl chloride, would have had further deleterious effects still on the atmospheric ozone shield. Upwelling of euxinic water may have released massive hydrogen sulphide emissions into the atmosphere and would poison terrestrial plants and animals and severely weaken the ozone layer, exposing much of the life that remained to fatal levels of UV radiation, although other modelling work has found that the release of this gas would not have significantly damaged the ozone layer. Indeed, biomarker evidence for anaerobic photosynthesis by Chlorobiaceae (green sulfur bacteria) from the Late-Permian into the Early Triassic indicates that hydrogen sulphide did upwell into shallow waters because these bacteria are restricted to the photic zone and use sulfide as an electron donor. Asteroid impact Evidence that an impact event may have caused the Cretaceous–Paleogene extinction has led to speculation that similar impacts may have been the cause of other extinction events, including the P–Tr extinction, and thus to a search for evidence of impacts at the times of other extinctions, such as large impact craters of the appropriate age. However, suggestions that an asteroid impact was the trigger of the Permian-Triassic extinction are now largely rejected. Reported evidence for an impact event from the P–Tr boundary level includes rare grains of shocked quartz in Australia and Antarctica; fullerenes trapping extraterrestrial noble gases; meteorite fragments in Antarctica; and grains rich in iron, nickel, and silicon, which may have been created by an impact. However, the accuracy of most of these claims has been challenged. For example, quartz from Graphite Peak in Antarctica, once considered "shocked", has been re-examined by optical and transmission electron microscopy. The observed features were concluded to be due not to shock, but rather to plastic deformation, consistent with formation in a tectonic environment such as volcanism. Iridium levels in many sites straddling the Permian-Triassic boundaries are not anomalous, providing evidence against an extraterrestrial impact as the PTME's cause. An impact crater on the seafloor would be evidence of a possible cause of the P–Tr extinction, but such a crater would by now have disappeared. As 70% of the Earth's surface is currently sea, an asteroid or comet fragment is now perhaps more than twice as likely to hit the ocean as it is to hit land. However, Earth's oldest ocean-floor crust is only 200 million years old as it is continually being destroyed and renewed by spreading and subduction. Furthermore, craters produced by very large impacts may be masked by extensive flood basalting from below after the crust is punctured or weakened. Yet, subduction should not be entirely accepted as an explanation for the lack of evidence: as with the K-T event, an ejecta blanket stratum rich in siderophilic elements (such as iridium) would be expected in formations from the time. A large impact might have triggered other mechanisms of extinction described above, such as the Siberian Traps eruptions at either an impact site or the antipode of an impact site. The abruptness of an impact also explains why more species did not rapidly evolve to survive, as would be expected if the Permian–Triassic event had been slower and less global than a meteorite impact. Bolide impact claims have been criticised on the grounds that they are unnecessary as explanations for the extinctions, and they do not fit the known data compatible with a protracted extinction spanning thousands of years. Additionally, many sites spanning the Permian-Triassic boundary display a complete lack of evidence of an impact event. Possible impact sites Possible impact craters proposed as the site of an impact causing the P–Tr extinction include the Bedout structure off the northwest coast of Australia and the hypothesized Wilkes Land crater of East Antarctica. An impact has not been proved in either case, and the idea has been widely criticized. The Wilkes Land geophysical feature is of very uncertain age, possibly later than the Permian–Triassic extinction. Another impact hypothesis postulates that the impact event which formed the Araguainha crater, whose formation has been dated to , a possible temporal range overlapping with the end-Permian extinction, precipitated the mass extinction. The impact occurred around extensive deposits of oil shale in the shallow marine Paraná–Karoo Basin, whose perturbation by the seismicity resulting from impact likely discharged about 1.6 teratonnes of methane into Earth's atmosphere, buttressing the already rapid warming caused by hydrocarbon release due to the Siberian Traps. The large earthquakes generated by the impact would have additionally generated massive tsunamis across much of the globe. Despite this, most palaeontologists reject the impact as being a significant driver of the extinction, citing the relatively low energy (equivalent to 105 to 106 of TNT, around two orders of magnitude lower than the impact energy believed to be required to induce mass extinctions) released by the impact. A 2017 paper noted the discovery of a circular gravity anomaly near the Falkland Islands which might correspond to an impact crater with a diameter of , as supported by seismic and magnetic evidence. Estimates for the age of the structure range up to 250 million years old. This would be substantially larger than the well-known Chicxulub impact crater associated with a later extinction. However, Dave McCarthy and colleagues from the British Geological Survey illustrated that the gravity anomaly is not circular and also that the seismic data presented by Rocca, Rampino and Baez Presser did not cross the proposed crater or provide any evidence for an impact crater. Methanogens A hypothesis published in 2014 posits that a genus of anaerobic methanogenic archaea known as Methanosarcina was responsible for the event. Three lines of evidence suggest that these microbes acquired a new metabolic pathway via gene transfer at about that time, enabling them to efficiently metabolize acetate into methane. That would have led to their exponential reproduction, allowing them to rapidly consume vast deposits of organic carbon that had accumulated in marine sediment. The result would have been a sharp buildup of methane and carbon dioxide in the oceans and atmosphere, in a manner that may be consistent with the 13C/12C isotopic record. Massive volcanism facilitated this process by releasing large amounts of nickel, a scarce metal which is a cofactor for enzymes involved in producing methane. Chemostratigraphic analysis of Permian-Triassic boundary sediments in Chaotian demonstrates a methanogenic burst could be responsible for some percentage of the carbon isotopic fluctuations. On the other hand, in the canonical Meishan sections, the nickel concentration increases somewhat after the concentrations have begun to fall. Interstellar dust John Gribbin argues that the Solar System last passed through a spiral arm of the Milky Way around 250 million years ago and that the resultant dusty gas clouds may have caused a dimming of the Sun, which combined with the effect of Pangaea to produce an ice age. Comparison to present global warming The PTME has been compared to the current anthropogenic global warming situation and Holocene extinction due to sharing the common characteristic of rapid rates of carbon dioxide release. Though the current rate of greenhouse gas emissions is more than an order of magnitude greater than the rate measured over the course of the PTME, the discharge of greenhouse gases during the PTME is poorly constrained geo-chronologically and was most likely pulsed and constrained to a few key, short intervals, rather than continuously occurring at a constant rate for the whole extinction interval; the rate of carbon release within these intervals was likely to have been similar in timing to modern anthropogenic emissions. As they did during the PTME, oceans in the present day are experiencing drops in pH and in oxygen levels, prompting further comparisons between modern anthropogenic ecological conditions and the PTME. Another biocalcification event similar in its effects on modern marine ecosystems is predicted to occur if carbon dioxide levels continue to rise. The changes in plant-insect interactions resulting from the PTME have also been invoked as possible indicators of the world's future ecology. The similarities between the two extinction events have led to warnings from geologists about the urgent need for reducing carbon dioxide emissions if an event similar to the PTME is to be prevented from occurring. Just as during the PTME, contemporary oceans experience their extreme change-change in the form of a decline in pH and oxygen levels, which further strengthens the pull between the two events. This is emphasised by geologist Lee Kump:“The Permian-Triassic mass extinction provides a stark reminder of the consequences of rapid carbon dioxide emissions. During the PTME, volcanic activity unleashed massive amounts of CO₂, leading to ocean acidification, deoxygenation, and widespread ecological collapse. Today, we see human activities driving similar processes at an even faster rate. The geological record shows that once these tipping points are reached, the cascading effects on ecosystems can last for millions of years.” If it continues to rise, the consequence could be another bio-calcification crisis, as seems to have occurred in the fossil record, which would have disastrous consequences for modern marine ecosystems.	Physical sciences	Geological history	null
24772	https://en.wikipedia.org/wiki/Phase%20modulation	Phase modulation	Phase modulation (PM) is a modulation pattern for conditioning communication signals for transmission. It encodes a message signal as variations in the instantaneous phase of a carrier wave. Phase modulation is one of the two principal forms of angle modulation, together with frequency modulation. In phase modulation, the instantaneous amplitude of the baseband signal modifies the phase of the carrier signal keeping its amplitude and frequency constant. The phase of a carrier signal is modulated to follow the changing signal level (amplitude) of the message signal. The peak amplitude and the frequency of the carrier signal are maintained constant, but as the amplitude of the message signal changes, the phase of the carrier changes correspondingly. Phase modulation is an integral part of many digital transmission coding schemes that underlie a wide range of technologies like Wi-Fi, GSM and satellite television. However it is not widely used for transmitting analog audio signals via radio waves. It is also used for signal and waveform generation in digital synthesizers, such as the Yamaha DX7, to implement FM synthesis. A related type of sound synthesis called phase distortion is used in the Casio CZ synthesizers. Foundation In general form, an analog modulation process of a sinusoidal carrier wave may be described by the following equation: . A(t) represents the time-varying amplitude of the sinusoidal carrier wave and the cosine-term is the carrier at its angular frequency , and the instantaneous phase deviation . This description directly provides the two major groups of modulation, amplitude modulation and angle modulation. In amplitude modulation, the angle term is held constant, while in angle modulation the term A(t) is constant and the second term of the equation has a functional relationship to the modulating message signal. The functional form of the cosine term, which contains the expression of the instantaneous phase as its argument, provides the distinction of the two types of angle modulation, frequency modulation (FM) and phase modulation (PM). In FM the message signal causes a functional variation of the carrier frequency. These variations are controlled by both the frequency and the amplitude of the modulating wave. In phase modulation, the instantaneous phase deviation (phase angle) of the carrier is controlled by the modulating waveform, such that the principal frequency remains constant. In principle, the modulating signal in both frequency and phase modulation may either be analog in nature, or it may be digital. The mathematics of the spectral behaviour reveals that there are two regions of particular interest: Modulation index As with other modulation indices, this quantity indicates by how much the modulated variable varies around its unmodulated level. It relates to the variations in the phase of the carrier signal: where is the peak phase deviation. Compare to the modulation index for frequency modulation.	Technology	Telecommunications	null
24776	https://en.wikipedia.org/wiki/Piston	Piston	A piston is a component of reciprocating engines, reciprocating pumps, gas compressors, hydraulic cylinders and pneumatic cylinders, among other similar mechanisms. It is the moving component that is contained by a cylinder and is made gas-tight by piston rings. In an engine, its purpose is to transfer force from expanding gas in the cylinder to the crankshaft via a piston rod and/or connecting rod. In a pump, the function is reversed and force is transferred from the crankshaft to the piston for the purpose of compressing or ejecting the fluid in the cylinder. In some engines, the piston also acts as a valve by covering and uncovering ports in the cylinder. Piston engines Internal combustion engines An internal combustion engine is acted upon by the pressure of the expanding combustion gases in the combustion chamber space at the top of the cylinder. This force then acts downwards through the connecting rod and onto the crankshaft. The connecting rod is attached to the piston by a swivelling gudgeon pin (US: wrist pin). This pin is mounted within the piston: unlike the steam engine, there is no piston rod or crosshead (except big two stroke engines). The typical piston design is on the picture. This type of piston is widely used in car diesel engines. According to purpose, supercharging level and working conditions of engines the shape and proportions can be changed. High-power diesel engines work in difficult conditions. Maximum pressure in the combustion chamber can reach 20 MPa and the maximum temperature of some piston surfaces can exceed 450 °C. It is possible to improve piston cooling by creating a special cooling cavity. Injector supplies this cooling cavity «A» with oil through oil supply channel «B». For better temperature reduction construction should be carefully calculated and analysed. Oil flow in the cooling cavity should be not less than 80% of the oil flow through the injector. The pin itself is of hardened steel and is fixed in the piston, but free to move in the connecting rod. A few designs use a 'fully floating' design that is loose in both components. All pins must be prevented from moving sideways and the ends of the pin digging into the cylinder wall, usually by circlips. Gas sealing is achieved by the use of piston rings. These are a number of narrow iron rings, fitted loosely into grooves in the piston, just below the crown. The rings are split at a point in the rim, allowing them to press against the cylinder with a light spring pressure. Two types of ring are used: the upper rings have solid faces and provide gas sealing; lower rings have narrow edges and a U-shaped profile, to act as oil scrapers. There are many proprietary and detail design features associated with piston rings. Pistons are usually cast or forged from aluminium alloys. For better strength and fatigue life, some racing pistons may be forged instead. Billet pistons are also used in racing engines because they do not rely on the size and architecture of available forgings, allowing for last-minute design changes. Although not commonly visible to the naked eye, pistons themselves are designed with a certain level of ovality and profile taper, meaning they are not perfectly round, and their diameter is larger near the bottom of the skirt than at the crown. Early pistons were of cast iron, but there were obvious benefits for engine balancing if a lighter alloy could be used. To produce pistons that could survive engine combustion temperatures, it was necessary to develop new alloys such as Y alloy and Hiduminium, specifically for use as pistons. A few early gas engines had double-acting cylinders, but otherwise effectively all internal combustion engine pistons are single-acting. During World War II, the US submarine Pompano was fitted with a prototype of the infamously unreliable H.O.R. double-acting two-stroke diesel engine. Although compact, for use in a cramped submarine, this design of engine was not repeated. Trunk pistons Trunk pistons are long relative to their diameter. They act both as a piston and cylindrical crosshead. As the connecting rod is angled for much of its rotation, there is also a side force that reacts along the side of the piston against the cylinder wall. A longer piston helps to support this. Trunk pistons have been a common design of piston since the early days of the reciprocating internal combustion engine. They were used for both petrol and diesel engines, although high speed engines have now adopted the lighter weight slipper piston. A characteristic of most trunk pistons, particularly for diesel engines, is that they have a groove for an oil ring below the gudgeon pin, in addition to the rings between the gudgeon pin and crown. The name 'trunk piston' derives from the 'trunk engine', an early design of marine steam engine. To make these more compact, they avoided the steam engine's usual piston rod with separate crosshead and were instead the first engine design to place the gudgeon pin directly within the piston. Otherwise these trunk engine pistons bore little resemblance to the trunk piston; they were extremely large diameter and double-acting. Their 'trunk' was a narrow cylinder mounted in the centre of the piston. Crosshead pistons Large slow-speed Diesel engines may require additional support for the side forces on the piston. These engines typically use crosshead pistons. The main piston has a large piston rod extending downwards from the piston to what is effectively a second smaller-diameter piston. The main piston is responsible for gas sealing and carries the piston rings. The smaller piston is purely a mechanical guide. It runs within a small cylinder as a trunk guide and also carries the gudgeon pin. Lubrication of the crosshead has advantages over the trunk piston as its lubricating oil is not subject to the heat of combustion: the oil is not contaminated by combustion soot particles, it does not break down owing to the heat and a thinner, less viscous oil may be used. The friction of both piston and crosshead may be only half of that for a trunk piston. Because of the additional weight of these pistons, they are not used for high-speed engines. Slipper pistons A slipper piston is a piston for a petrol engine that has been reduced in size and weight as much as possible. In the extreme case, they are reduced to the piston crown, support for the piston rings, and just enough of the piston skirt remaining to leave two lands so as to stop the piston rocking in the bore. The sides of the piston skirt around the gudgeon pin are reduced away from the cylinder wall. The purpose is mostly to reduce the reciprocating mass, thus making it easier to balance the engine and so permit high speeds. In racing applications, slipper piston skirts can be configured to yield extremely light weight while maintaining the rigidity and strength of a full skirt. Reduced inertia also improves mechanical efficiency of the engine: the forces required to accelerate and decelerate the reciprocating parts cause more piston friction with the cylinder wall than the fluid pressure on the piston head. A secondary benefit may be some reduction in friction with the cylinder wall, since the area of the skirt, which slides up and down in the cylinder is reduced by half. However, most friction is due to the piston rings, which are the parts which actually fit the tightest in the bore and the bearing surfaces of the wrist pin, and thus the benefit is reduced. Deflector pistons Deflector pistons are used in two-stroke engines with crankcase compression, where the gas flow within the cylinder must be carefully directed in order to provide efficient scavenging. With cross scavenging, the transfer (inlet to the cylinder) and exhaust ports are on directly facing sides of the cylinder wall. To prevent the incoming mixture passing straight across from one port to the other, the piston has a raised rib on its crown. This is intended to deflect the incoming mixture upwards, around the combustion chamber. Much effort, and many different designs of piston crown, went into developing improved scavenging. The crowns developed from a simple rib to a large asymmetric bulge, usually with a steep face on the inlet side and a gentle curve on the exhaust. Despite this, cross scavenging was never as effective as hoped. Most engines today use Schnuerle porting instead. This places a pair of transfer ports in the sides of the cylinder and encourages gas flow to rotate around a vertical axis, rather than a horizontal axis. Racing pistons In racing engines, piston strength and stiffness is typically much higher than that of a passenger car engine, while the weight is much less, to achieve the high engine RPM necessary in racing. Hydraulic cylinders Hydraulic cylinders can be both single-acting or double-acting. A hydraulic actuator controls the movement of the piston back and/or forth. Guide rings guides the piston and rod and absorb the radial forces that act perpendicularly to the cylinder and prevent contact between sliding the metal parts. Steam engines Steam engines are usually double-acting (i.e. steam pressure acts alternately on each side of the piston) and the admission and release of steam is controlled by slide valves, piston valves or poppet valves. Consequently, steam engine pistons are nearly always comparatively thin discs: their diameter is several times their thickness. (One exception is the trunk engine piston, shaped more like those in a modern internal-combustion engine.) Another factor is that since almost all steam engines use crossheads to translate the force to the drive rod, there are few lateral forces acting to try and "rock" the piston, so a cylinder-shaped piston skirt isn't necessary. Pumps Piston pumps can be used to move liquids or compress gases. For liquids For gases Air cannons There are two special type of pistons used in air cannons: close tolerance pistons and double pistons. In close tolerance pistons O-rings serve as a valve, but O-rings are not used in double piston types.	Technology	Rigid components	null
24823	https://en.wikipedia.org/wiki/Pterodactylus	Pterodactylus	Pterodactylus (from ) is a genus of extinct pterosaurs. It is thought to contain only a single species, Pterodactylus antiquus, which was the first pterosaur to be named and identified as a flying reptile and one of the first prehistoric reptiles to ever be discovered. Fossil remains of Pterodactylus have primarily been found in the Solnhofen limestone of Bavaria, Germany, which dates from the Late Jurassic period (Tithonian stage), about 150.8 to 148.5 million years ago. More fragmentary remains of Pterodactylus have tentatively been identified from elsewhere in Europe and in Africa. Pterodactylus was a generalist carnivore that probably fed on a variety of invertebrates and vertebrates. Like all pterosaurs, Pterodactylus had wings formed by a skin and muscle membrane stretching from its elongated fourth finger to its hind limbs. It was supported internally by collagen fibres and externally by keratinous ridges. Pterodactylus was a small pterosaur compared to other famous genera such as Pteranodon and Quetzalcoatlus, and it also lived earlier, during the Late Jurassic period, while both Pteranodon and Quetzalcoatlus lived during the Late Cretaceous. Pterodactylus lived alongside other small pterosaurs such as the well-known Rhamphorhynchus, as well as other genera such as Scaphognathus, Anurognathus and Ctenochasma. Pterodactylus is classified as an early-branching member of the ctenochasmatid lineage, within the pterosaur clade Pterodactyloidea. Discovery and history The type specimen of the animal now known as Pterodactylus antiquus was the first pterosaur fossil ever to be identified. The first Pterodactylus specimen was described by the Italian scientist in 1784, based on a fossil skeleton that had been unearthed from the Solnhofen limestone of Bavaria. Collini was the curator of the , or nature cabinet of curiosities (a precursor to the modern concept of the natural history museum), in the palace of Charles Theodore, Elector of Bavaria at Mannheim. The specimen had been given to the collection by Count around 1780, having been recovered from a lithographic limestone quarry in . The actual date of the specimen's discovery and entry into the collection is unknown however, and it was not mentioned in a catalogue of the collection taken in 1767, so it must have been acquired at some point between that date and its 1784 description by Collini. This makes it potentially the earliest documented pterosaur find; the "Pester Exemplar" of the genus Aurorazhdarcho was described in 1779 and possibly discovered earlier than the Mannheim specimen, but it was at first considered to be a fossilized crustacean, and it was not until 1856 that this species was properly described as a pterosaur by German paleontologist . In his first description of the Mannheim specimen, Collini did not conclude that it was a flying animal. In fact, Collini could not fathom what kind of animal it might have been, rejecting affinities with the birds or the bats. He speculated that it may have been a sea creature, not for any anatomical reason, but because he thought the ocean depths were more likely to have housed unknown types of animals. The idea that pterosaurs were aquatic animals persisted among a minority of scientists as late as 1830, when the German zoologist Johann Georg Wagler published a text on "amphibians" which included an illustration of Pterodactylus using its wings as flippers. Wagler went so far as to classify Pterodactylus, along with other aquatic vertebrates (namely plesiosaurs, ichthyosaurs, and monotremes), in the class Gryphi, between birds and mammals. The German/French scientist Johann Hermann was the one who first stated that Pterodactylus used its long fourth finger to support a wing membrane. Back in March 1800, Hermann alerted the prominent French scientist Georges Cuvier to the existence of Collini's fossil, believing that it had been captured by the invading forces of the French Consulate and sent to collections in Paris (and perhaps to Cuvier himself) as war booty; at the time special French political commissars systematically seized art treasures and objects of scientific interest. Hermann sent Cuvier a letter containing his own interpretation of the specimen (though he had not examined it personally), which he believed to be a mammal, including the first known life restoration of a pterosaur. Hermann restored the animal with wing membranes extending from the long fourth finger to the ankle and a covering of fur (neither wing membranes nor fur had been preserved in the specimen). Hermann also added a membrane between the neck and wrist, as is the condition in bats. Cuvier agreed with this interpretation, and at Hermann's suggestion, Cuvier became the first to publish these ideas in December 1800 in a very short description. However, contrary to Hermann, Cuvier was convinced the animal was a reptile. The specimen had not in fact been seized by the French. Rather, in 1802, following the death of Charles Theodore, it was brought to Munich, where Baron Johann Paul Carl von Moll had obtained a general exemption of confiscation for the Bavarian collections. Cuvier asked von Moll to study the fossil but was informed it could not be found. In 1809 Cuvier published a somewhat longer description, in which he named the animal Petro-Dactyle, this was a typographical error however, and was later corrected by him to Ptéro-Dactyle. He also refuted a hypothesis by Johann Friedrich Blumenbach that it would have been a shore bird. Cuvier remarked: "It is not possible to doubt that the long finger served to support a membrane that, by lengthening the anterior extremity of this animal, formed a good wing." Contrary to von Moll's report, the fossil was not missing; it was being studied by Samuel Thomas von Sömmerring, who gave a public lecture about it on December 27, 1810. In January 1811, von Sömmerring wrote a letter to Cuvier deploring the fact that he had only recently been informed of Cuvier's request for information. His lecture was published in 1812, and in it von Sömmerring named the species Ornithocephalus antiquus. The animal was described as being both a bat, and a form in between mammals and birds, i.e. not intermediate in descent but in "affinity" or archetype. Cuvier disagreed, and the same year in his Ossemens fossiles provided a lengthy description in which he restated that the animal was a reptile. It was not until 1817 that a second specimen of Pterodactylus came to light, again from Solnhofen. This tiny specimen was that year described by von Sömmerring as Ornithocephalus brevirostris, named for its short snout, now understood to be a juvenile character (this specimen is now thought to represent a juvenile specimen of a different genus, probably Ctenochasma). He provided a restoration of the skeleton, the first one published for any pterosaur. This restoration was very inaccurate, von Sömmerring mistaking the long metacarpals for the bones of the lower arm, the lower arm for the humerus, this upper arm for the breast bone and this sternum again for the shoulder blades. Sömmerring did not change his opinion that these forms were bats and this "bat model" for interpreting pterosaurs would remain influential long after a consensus had been reached around 1860 that they were reptiles. The standard assumptions were that pterosaurs were quadrupedal, clumsy on the ground, furred, warmblooded and had a wing membrane reaching the ankle. Some of these elements have been confirmed, some refuted by modern research, while others remain disputed. In 1815, the generic name Ptéro-Dactyle was latinized to Pterodactylus by Constantine Samuel Rafinesque. Unaware of Rafinesque's publication however, Cuvier himself in 1819 latinized the name Ptéro-Dactyle again to Pterodactylus, but the specific name he then gave, longirostris, has to give precedence to von Sömmerring's antiquus. In 1888, English naturalist Richard Lydekker designated Pterodactylus antiquus as the type species of Pterodactylus, and considered Ornithocephalus antiquus a synonym. He also designated specimen BSP AS.I.739 as the holotype of the genus. Description Pterodactylus is known from over 30 fossil specimens, and though most belong to juveniles, many preserve complete skeletons. Pterodactylus antiquus was a relatively small pterosaur, with an estimated adult wingspan of about , based on the only known adult specimen, which is represented by an isolated skull. Other "species" were once thought to have been smaller. However, these smaller specimens have been shown to represent juveniles of Pterodactylus, as well as its contemporary relatives including Ctenochasma, Germanodactylus, Aurorazhdarcho, Gnathosaurus, and hypothetically Aerodactylus if this genus is truly valid. The skulls of adult Pterodactylus were long and thin, with about 90 narrow and conical teeth. The teeth extended back from the tips of both jaws, and became smaller farther away from the jaw tips. This was unlike the ones seen in most relatives, where teeth were absent in the upper jaw tip and were relatively uniform in size. The teeth of Pterodactylus also extended farther back into the jaw compared to close relatives, and some were present below the front of the nasoantorbital fenestra, which is the largest opening in the skull. Another autapomorphy that Pterodactylus has is that the skull and jaws were straight, which are unlike the upwardly curved jaws seen in the related ctenochasmatids. Pterodactylus, like related pterosaurs, had a crest on its skull composed mainly of soft tissues. In adult Pterodactylus, this crest extended between the back edge of the antorbital fenestra and the back of the skull. In at least one specimen, the crest had a short bony base, also seen in related pterosaurs like Germanodactylus. Solid crests have only been found on large, fully adult specimens of Pterodactylus, indicating that this was a display structure that became larger and more well developed as individuals reached maturity. In 2013, pterosaur researcher S. Christopher Bennett noted that other authors claimed that the soft tissue crest of Pterodactylus extended backward behind the skull; Bennett himself, however, didn't find any evidence for the crest extending past the back of the skull. Two specimens of P. antiquus (the holotype specimen BSP AS I 739 and the incomplete skull BMMS 7, the largest known skull of P. antiquus) have a low bony crest on their skulls; in BMMS 7 it is 47.5 mm long (1.87 inches, more or less 24% of the estimated total length of its skull) and has a maximum height of 0.9 mm (0.035 inches) above the orbit. Several specimens previously referred to P. antiquus preserved evidence of the soft tissue extensions of these crests, including an "occipital lappet", a flexible, tab-like structure extending from the back of the skull. Most of these specimens have been reclassified in the related species Aerodactylus scolopaciceps, which may however be nothing more than a junior synonym. Even if Aerodactylus were valid, at least one specimen with these features is still considered to belong to Pterodactylus, BSP 1929 I 18, which has an occipital lappet similar to the proposed Aerodactylus definition, and also possesses a small triangular soft tissue crest with the peak of the crest positioned above the eyes. Paleobiology Life history Like other pterosaurs (most notably Rhamphorhynchus), Pterodactylus specimens can vary considerably based on age or level of maturity. Both the proportions of the limb bones, size and shape of the skull, and size and number of teeth changed as the animals grew. Historically, this has led to various growth stages (including growth stages of related pterosaurs) being mistaken for new species of Pterodactylus. Several detailed studies using various methods to measure growth curves among known specimens have suggested that there is actually only one valid species of Pterodactylus, P. antiquus. The youngest immature specimens of Pterodactylus antiquus (alternately interpreted as young specimens of the distinct species P. kochi) have a small number of teeth, as few as 15 in some, and the teeth have a relatively broad base. The teeth of other P. antiquus specimens are both narrower and more numerous (up to 90 teeth are present in several specimens). Pterodactylus specimens can be divided into two distinct year classes. In the first year class, the skulls are only in length. The second year class is characterized by skulls of around long, but are still immature however. These first two size groups were once classified as juveniles and adults of the species P. kochi, until further study showed that even the supposed "adults" were immature, and possibly belong to a distinct genus. A third year class is represented by specimens of the "traditional" P. antiquus, as well as a few isolated, large specimens once assigned to P. kochi that overlap P. antiquus in size. However, all specimens in this third year class also show sign of immaturity. Fully mature Pterodactylus specimens remain unknown, or may have been mistakenly classified as a different genus. Growth and breeding seasons The distinct year classes of Pterodactylus antiquus specimens show that this species, like the contemporary Rhamphorhynchus muensteri, likely bred seasonally and grew consistently during its lifetime. A new generation of 1st year class P. antiquus would have been produced seasonally, and reached 2nd-year size by the time the next generation hatched, creating distinct 'clumps' of similarly-sized and aged individuals in the fossil record. The smallest size class probably consisted of individuals that had just begun to fly and were less than one year old. The second year class represents individuals one to two years old, and the rare third year class is composed of specimens over two years old. This growth pattern is similar to modern crocodilians, rather than the rapid growth of modern birds. Daily activity patterns Comparisons between the scleral rings of Pterodactylus antiquus and modern birds and reptiles suggest that it may have been diurnal. This may also indicate niche partitioning with contemporary pterosaurs inferred to be nocturnal, such as Ctenochasma and Rhamphorhynchus. Diet Based on the shape, size, and arrangement of its teeth, Pterodactylus has long been recognized as a carnivore specializing in small animals. A 2020 study of pterosaur tooth wear supported the hypothesis that Pterodactylus preyed mainly on invertebrates and had a generalist feeding strategy, indicated by a relatively high bite force. Paleoecology Specimens of Pterodactylus have been found mainly in the Solnhofen limestone (geologically known as the Altmühltal Formation) of Bavaria, Germany. The main composition of this formation is fine-grained limestone that originated mainly from the nearby towns Solnhofen and Eichstätt, which is formed by mud silt deposits. The Solnhofen Limestone is a diverse Lagerstätte that contains a wide range of different creatures, including highly detailed fossilized imprints of soft bodied organisms such as jellyfishes. Abundant specimens of pterosaurs similar to Pterodactylus were also found within the formation, these include the rhamphorhynchids Rhamphorhynchus and Scaphognathus, several gallodactylids such as Aerodactylus, Ardeadactylus, Aurorazhdarcho and Cycnorhamphus, the ctenochasmatids Ctenochasma and Gnathosaurus, the anurognathid Anurognathus, the germanodactylid Germanodactylus, as well as the basal euctenochasmatian Diopecephalus. Fossil remains of the dinosaurs Archaeopteryx and Compsognathus were also found within the limestone, these specimens were related to early evolution of feathers, since they were some of the only ones that had them during the Jurassic period. Various lizard remains were also found alongside those of Pterodactylus, with several specimens assigned to Ardeosaurus, Bavarisaurus and Eichstaettisaurus. Crocodylomorph specimens were widely distributed within the fossil site, most were assigned to the metriorhynchid genera Cricosaurus, Dakosaurus, Geosaurus and Rhacheosaurus. These genera are colloquially called as marine or sea crocodiles due to their similar built. The turtle genera Eurysternum and Paleomedusa were also found within the formation. Fossils of the ichthyosaur Aegirosaurus also appeared to be present in the site, as well as fish remains, with many specimens assigned to ray-finned fishes such as the halecomorphs Lepidotes, Propterus, Gyrodus, Mesturus, Proscinetes, Caturus, Ophiopsis and Ophiopsiella, the pachycormids Asthenocormus, Hypsocormus and Orthocormus, as well as the aspidorhynchid Aspidorhynchus, and the ichthyodectid Thrissops. Classification Initial classifications for Pterodactylus started when paleontologist Hermann von Meyer used the name Pterodactyli to contain Pterodactylus and other pterosaurs known at the time. This was emended to the family Pterodactylidae by Prince Charles Lucien Bonaparte in 1838. However, this group has more recently been given several competing definitions. Beginning in 2014, researchers Steven Vidovic and David Martill constructed an analysis in which several pterosaurs traditionally thought of as archaeopterodactyloids closely related to the ctenochasmatoids may have been more closely related to the more advanced dsungaripteroids, or in some cases, fall outside both groups. Their conclusion was published in 2017, in which they placed Pterodactylus as a basal member of the suborder Pterodactyloidea. As illustrated below, the results of a different topology are based on a phylogenetic analysis made by Longrich, Martill, and Andres in 2018. Unlike the previous results above, they placed Pterodactylus within the clade Euctenochasmatia, resulting in a more derived position. Formerly assigned species Numerous species have been assigned to Pterodactylus in the years since its discovery. In the first half of the 19th century any new pterosaur species would be named Pterodactylus, which thus became a "wastebasket taxon". Even after clearly different forms had later been given their own generic name, new species would be created from the very productive sites, throughout Europe and North America, often based on only slightly different material. The earliest reassignments of pterosaur species to Pterodactylus started in 1825, with the description of Rhamphorhynchus; fossil collector Georg Graf zu Münster alerted the German paleontologist Samuel Thomas von Sömmerring about several distinct fossil specimens, Sömmerring thought that they belonged to an ancient bird. Further fossil preparations had uncovered teeth, to which Graf zu Münster created a skull cast. He later sent the cast to Professor Georg August Goldfuss, who recognized it as a pterosaur, specifically a species of Pterodactylus. At the time however, most paleontologists incorrectly consider the genus Ornithocephalus () to be the valid name for Pterodactylus, and therefore the specimen found was named as Ornithocephalus Münsteri, which was first mentioned by Graf zu Münster himself. Another specimen was found and described by Graf zu Münster in 1839, he assigned this specimen to a new separate species called Ornithocephalus longicaudus; the specific name means 'long tail', in reference to the animal's tail size. German paleontologist Hermann von Meyer in 1845 officially emended that the genus Pterodactylus had priority over Ornithocephalus, so he reassigned the species O. münsteri and O. longicaudus into Pterodactylus münsteri and Pterodactylus longicaudus. In 1846, von Meyer created the new species Pterodactylus gemmingi based on long-tailed remains; the specific name honors the fossil collector Carl Eming von Gemming. Later, in 1847, von Meyer finally erected the generic name Rhamphorhynchus () due to the distinctively long tails seen in the specimens found, which are much longer than those seen in Pterodactylus. He assigned the species P. longicaudus as the type species of Rhamphorhynchus, which resulted in a new combination called Rhamphorhynchus longicaudus. The species R. münsteri was later changed to R. muensteri by Lydekker in 1888, due to the ICZN rule that prohibits non-standard Latin characters, such as ü, in scientific names. Beginning in 1846, many pterosaur specimens were found near the village of Burham in Kent, England by British paleontologists James Scott Bowerbank and Sir Richard Owen. Bowerbank had assigned fossil remains to two new species; the first was named in 1846 as Pterodactylus giganteus; the specific name means 'the gigantic one' in Latin, in reference to the large size of the remains, and the second species was named in 1851 as Pterodactylus cuvieri, in honor of the French scientist Georges Cuvier. Later in 1851, Owen named and described new pterosaur specimens that have been found yet again in England. He assigned these specimens to a new species called Pterodactylus compressirostris. In 1914 however, paleontologist Reginald Hooley redescribed P. compressirostris, to which he erected the genus Lonchodectes (), and therefore made P. compressirostris the type species, and created the new combination L. compressirostris. In a 2013 review, P. giganteus and P. cuvieri were reassigned to new genera; P. giganteus was reassigned to a genus called Lonchodraco ('lance dragon'), which resulted in a new combination called L. giganteus, and P. cuvieri was reassigned to the new genus Cimoliopterus ('chalk wing'), creating C. cuvieri. Back in 1859, Owen had found remains the front part of a snout in the Cambridge Greensand, and assigned it into the species Pterodactylus segwickii; in honor of Adam Sedgwick, a British geologist. This species however, was reassigned to the genus Camposipterus in 2013, therefore creating the new combination Camposipterus segwickii. Later, in 1861, Owen had uncovered multiple distinctively looking fossil remains yet again in the Cambridge Greensand, these were assigned to a new species named Pterodactylus simus, though the British paleontologist Harry Govier Seeley had created a separate generic name called Ornithocheirus, and reassigned P. simus as the type species, which created the combination Ornithocheirus simus. Between the years 1869 and 1870, Seeley had reassigned many pterosaur species into Ornithocheirus, while also creating several new species. Many of these species however, are now reclassified to other genera, or considered . In 1874, further specimens were found in England, again by Owen, these ones were assigned to a new species called Pterodactylus sagittirostris, this species however, was reassigned to the genus Lonchodectes in 1914 by Hooley, which resulted in an L. sagittirostris. This conclusion was revised by Rigal et al. in 2017, who disagreed with Hooley's reassignment, and therefore created the genus Serradraco, which afterwards resulted in a new combination called S. sagittirostris. Assigning new pterosaur species to Pterodactylus was not only common in Europe, but also in North America; paleontologists such as Othniel Charles Marsh in 1871 for example, described several toothless pterosaur specimens, which were accompanied by teeth that belonged to the fish Xiphactinus, which Marsh assumed that these teeth belonged to the pterosaur specimens he found, since all pterosaurs discovered at the time had teeth. He then assigned these specimens to a new species called "Pterodactylus oweni", but this was changed to Pterodactylus occidentalis because "P. oweni" was found to have been preoccupied by a pterosaur species described with the same name back in 1864 by Seeley. In 1872, American paleontologist Edward Drinker Cope also found various pterosaur specimens in North America, he assigned these to two new species known as Ornithochirus umbrosus and Ornithochirus harpyia, Cope attempted to assign the specimens he found to the genus Ornithocheirus, but misspelled forgetting the 'e'. In 1875 however, Cope reassigned the species O. umbrosus and O. harpyia into Pterodactylus umbrosus and Pterodactylus harpyia, though these species had been considered ever since. Paleontologist Samuel Wendell Williston unearthed the first skull of the pterosaur, and found that the animal was toothless, this made Marsh create the genus Pteranodon (lit. 'toothless wing'), and therefore reassigned all the American pterosaur species, including the ones that he named, from Pterodactylus to Pteranodon. Later, in the 1980s, subsequent revisions by Peter Wellnhofer had reduced the number of recognized species to about half a dozen. Many species assigned to Pterodactylus had been based on juvenile specimens, and subsequently been recognized as immature individuals of other species or genera. By the 1990s it was understood that this was even true for part of the remaining species. P. elegans, for example, was found by numerous studies to be an immature Ctenochasma. Another species of Pterodactylus originally based on small, immature specimens was P. micronyx. However, it has been difficult to determine exactly of what genus and species P. micronyx might be the juvenile form. Stéphane Jouve, Christopher Bennett and others had once suggested that it probably belonged either to Gnathosaurus subulatus or one of the species belonging to Ctenochasma, though after additional research Bennett assigned it to the genus Aurorazhdarcho. Another species with a complex history is P. longicollum, named by von Meyer in 1854, based on a large specimen with a long neck and fewer teeth. Many researchers, including David Unwin, have found P. longicollum to be distinct from P. kochi and P. antiquus. Unwin found P. longicollum to be closer to Germanodactylus and therefore requiring a new genus name. It has sometimes been placed in the genus Diopecephalus because Harry Govier Seeley based this genus partly on the P. longicollum material. However, it was shown by Bennett that the type specimen later designated for Diopecephalus was a fossil belonging to P. kochi, and no longer thought to be separate from Pterodactylus. Diopecephalus is therefore a synonym of Pterodactylus, and as such is unavailable for use as a new genus for "P." longicollum. "P." longicollum was eventually made the type species of a separate genus Ardeadactylus. Controversial species The only well-known and well-supported species left by the first decades of the 21st century were P. antiquus and P. kochi. However, most studies between 1995 and 2010 found little reason to separate even these two species, and treated them as synonymous. More recent studies of pterosaur relationships have found anurognathids and pterodactyloids to be sister groups, which would limit the more inclusive group Caelidracones to just two clades. In 1996, Bennett suggested that the differences between specimens of P. kochi and P. antiquus could be explained by differences in age, with P. kochi (including specimens alternately classified in the species P. scolopaciceps) representing an immature growth stage of P. antiquus. In a 2004 paper, Jouve used a different method of analysis and recovered the same result, showing that the "distinctive" features of P. kochi were age-related, and using mathematical comparison to show that the two forms are different growth stages of the same species. An additional review of the specimens published in 2013 demonstrated that some of the supposed differences between P. kochi and P. antiquus were due to measurement errors, further supporting their synonymy. By the 2010s, a large body of research had been developed based on the idea that P. kochi and P. scolopaciceps were early growth stages of P. antiquus. However, in 2014, two scientists began publishing research that challenged this paradigm. Steven Vidovic and David Martill concluded that differences between specimens of P. kochi, P. scolopaciceps, and P. antiquus, such as different lengths of neck vertebrae, thinner or thicker teeth, more rounded skulls, and how far the teeth extended back in the jaws, were significant enough to separate them into three distinct species. Vidovic and Martill also performed a phylogenetic analysis which treated all relevant specimens as distinct units, and found that the P. kochi type specimen did not form a natural group with that of P. antiquus. They concluded that the genus Diopecephalus could be returned to use to distinguish "P". kochi from P. antiquus. They named the new genus Aerodactylus for P. scolopaciceps as well. So, what Bennett considered early growth stages of one species, Vidovic and Martill considered representatives of new species. In 2017, Bennett challenged this hypothesis, he claimed that while Vidovic and Martill had identified real differences between these three groups of specimens, they had not provided any rationale that the differences were enough to distinguish them as species, rather than just individual variation, growth changes, or simply due to crushing and distortion during the fossilization process. Bennett pointed in particular to the data used to distinguish Aerodactylus, which was so different from the data for related species, it might be due to an unnatural assemblage of specimens. As a result, Bennett continued to consider Diopecephalus and Aerodactylus simply as year-classes of immature Pterodactylus antiquus. List of species During its over-200-year history, the various species of Pterodactylus have gone through a number of changes in classification and thus have acquired a large number of synonyms. Additionally, a number of species assigned to Pterodactylus are based on poor remains that have proven difficult to assign to one species or another and are therefore considered (). The following list includes names that were used to identify new pterosaur species that now have been reclassified, or until recently thought to be pertaining to Pterodactylus proper, and names based on other material that has as yet not been assigned to other genera. This list also includes species that are ('naked names'), which are species that were not published formally. Species that are ('forgotten names') are the ones that have been disused, and species that are ('rejected names') are the ones that have been rejected because a more preferable name had been accepted instead. {\| class="wikitable sortable" \|- ! Name ! Author ! Year ! Status ! class="unsortable" \|	Biology and health sciences	Pterosaurs	Animals
24824	https://en.wikipedia.org/wiki/Pterosaur	Pterosaur	Pterosaurs are an extinct clade of flying reptiles in the order Pterosauria. They existed during most of the Mesozoic: from the Late Triassic to the end of the Cretaceous (228 million to 66 million years ago). Pterosaurs are the earliest vertebrates known to have evolved powered flight. Their wings were formed by a membrane of skin, muscle, and other tissues stretching from the ankles to a dramatically lengthened fourth finger. There were two major types of pterosaurs. Basal pterosaurs (also called 'non-pterodactyloid pterosaurs' or 'rhamphorhynchoids') were smaller animals with fully toothed jaws and, typically, long tails. Their wide wing membranes probably included and connected the hind legs. On the ground, they would have had an awkward sprawling posture, but the anatomy of their joints and strong claws would have made them effective climbers, and some may have even lived in trees. Basal pterosaurs were insectivores or predators of small vertebrates. Later pterosaurs (pterodactyloids) evolved many sizes, shapes, and lifestyles. Pterodactyloids had narrower wings with free hind limbs, highly reduced tails, and long necks with large heads. On the ground, they walked well on all four limbs with an upright posture, standing plantigrade on the hind feet and folding the wing finger upward to walk on the three-fingered "hand". They could take off from the ground, and fossil trackways show that at least some species were able to run, wade, and/or swim. Their jaws had horny beaks, and some groups lacked teeth. Some groups developed elaborate head crests with sexual dimorphism. Pterosaurs sported coats of hair-like filaments known as pycnofibers, which covered their bodies and parts of their wings. Pycnofibers grew in several forms, from simple filaments to branching down feathers. These may be homologous to the down feathers found on both avian and some non-avian dinosaurs, suggesting that early feathers evolved in the common ancestor of pterosaurs and dinosaurs, possibly as insulation. They were warm-blooded (endothermic), active animals. The respiratory system had efficient unidirectional "flow-through" breathing using air sacs, which hollowed out their bones to an extreme extent. Pterosaurs spanned a wide range of adult sizes, from the very small anurognathids to the largest known flying creatures, including Quetzalcoatlus and Hatzegopteryx, which reached wingspans of at least nine metres. The combination of endothermy, a good oxygen supply and strong muscles made pterosaurs powerful and capable flyers. Pterosaurs are often referred to by popular media or the general public as "flying dinosaurs", but dinosaurs are defined as the descendants of the last common ancestor of the Saurischia and Ornithischia, which excludes the pterosaurs. Pterosaurs are nonetheless more closely related to birds and other dinosaurs than to crocodiles or any other living reptile, though they are not bird ancestors. Pterosaurs are also colloquially referred to as pterodactyls, particularly in fiction and journalism. However, technically, pterodactyl may refer to members of the genus Pterodactylus, and more broadly to members of the suborder Pterodactyloidea of the pterosaurs. Pterosaurs had a variety of lifestyles. Traditionally seen as fish-eaters, the group is now understood to have also included hunters of land animals, insectivores, fruit eaters and even predators of other pterosaurs. They reproduced by eggs, some fossils of which have been discovered. Description The anatomy of pterosaurs was highly modified from their reptilian ancestors by the adaptation to flight. Pterosaur bones were hollow and air-filled, like those of birds. This provided a higher muscle attachment surface for a given skeletal weight. The bone walls were often paper-thin. They had a large and keeled breastbone for flight muscles and an enlarged brain able to coordinate complex flying behaviour. Pterosaur skeletons often show considerable fusion. In the skull, the sutures between elements disappeared. In some later pterosaurs, the backbone over the shoulders fused into a structure known as a notarium, which served to stiffen the torso during flight, and provide a stable support for the shoulder blade. Likewise, the sacral vertebrae could form a single synsacrum while the pelvic bones fused also. Basal pterosaurs include the clades Dimorphodontidae (Dimorphodon), Campylognathididae (Eudimorphodon, Campyognathoides), and Rhamphorhynchidae (Rhamphorhynchus, Scaphognathus). Pterodactyloids include the clades Ornithocheiroidea (Istiodactylus, Ornithocheirus, Pteranodon), Ctenochasmatoidea (Ctenochasma, Pterodactylus), Dsungaripteroidea (Germanodactylus, Dsungaripterus), and Azhdarchoidea (Tapejara, Tupuxuara, Quetzalcoatlus). The two groups overlapped in time, but the earliest pterosaurs in the fossil record are basal pterosaurs, and the latest pterosaurs are pterodactyloids. The position of the clade Anurognathidae (Anurognathus, Jeholopterus, Vesperopterylus) is debated. Anurognathids were highly specialized. Small flyers with shortened jaws and a wide gape, some had large eyes suggesting nocturnal or crepuscular habits, mouth bristles, and feet adapted for clinging. Parallel adaptations are seen in birds and bats that prey on insects in flight. Size Pterosaurs had a wide range of sizes, though they were generally large. The smallest species had a wingspan no less than . The most sizeable forms represent the largest known animals ever to fly, with wingspans of up to . Standing, such giants could reach the height of a modern giraffe. Traditionally, it was assumed that pterosaurs were extremely light relative to their size. Later, it was understood that this would imply unrealistically low densities of their soft tissues. Some modern estimates therefore extrapolate a weight of up to for the largest species. Skull, teeth, and crests Compared to the other vertebrate flying groups, the birds and bats, pterosaur skulls were typically quite large. Most pterosaur skulls had elongated jaws. Their skull bones tend to be fused in adult individuals. Early pterosaurs often had heterodont teeth, varying in build, and some still had teeth in the palate. In later groups the teeth mostly became conical. Front teeth were often longer, forming a "prey grab" in transversely expanded jaw tips, but size and position were very variable among species. With the derived Pterodactyloidea, the skulls became even more elongated, sometimes surpassing the combined neck and torso in length. This was caused by a stretching and fusion of the front snout bone, the premaxilla, with the upper jawbone, the maxilla. Unlike most archosaurs, the nasal and antorbital openings of pterodactyloid pterosaurs merged into a single large opening, called the nasoantorbital fenestra. This feature likely evolved to lighten the skull for flight. In contrast, the bones behind the eye socket contracted and rotated, strongly inclining the rear skull and bringing the jaw joint forward. The braincase was relatively large for reptiles. In some cases, fossilized keratinous beak tissue has been preserved, though in toothed forms, the beak is small and restricted to the jaw tips and does not involve the teeth. Some advanced beaked forms were toothless, such as the Pteranodontidae and Azhdarchidae, and had larger, more extensive, and more bird-like beaks. Some groups had specialised tooth forms. The Istiodactylidae had recurved teeth for eating meat. Ctenochasmatidae used combs of numerous needle-like teeth for filter feeding; Pterodaustro could have over a thousand bristle-like teeth. Dsungaripteridae covered their teeth with jawbone tissue for a crushing function. If teeth were present, they were placed in separate tooth sockets. Replacement teeth were generated behind, not below, the older teeth. The public image of pterosaurs is defined by their elaborate head crests. This was influenced by the distinctive backward-pointing crest of the well-known Pteranodon. The main positions of such crests are the front of the snout, as an outgrowth of the premaxillae, or the rear of the skull as an extension of the parietal bones in which case it is called a "supraoccipital crest". Front and rear crests can be present simultaneously and might be fused into a single larger structure, the most expansive of which is shown by the Tapejaridae. Nyctosaurus sported a bizarre antler-like crest. The crests were only a few millimetres thin transversely. The bony crest base would typically be extended by keratinous or other soft tissue. Since the 1990s, new discoveries and a more thorough study of old specimens have shown that crests are far more widespread among pterosaurs than previously assumed. That they were extended by or composed completely of keratin, which does not fossilize easily, had misled earlier research. For Pterorhynchus and Pterodactylus, the true extent of these crests has only been uncovered using ultraviolet photography. While fossil crests used to be restricted to the more advanced Pterodactyloidea, Pterorhynchus and Austriadactylus show that even some early pterosaurs possessed them. Like the upper jaws, the paired lower jaws of pterosaurs were very elongated. In advanced forms, they tended to be shorter than the upper cranium because the jaw joint was in a more forward position. The front lower jaw bones, the dentaries or ossa dentalia, were at the tip tightly fused into a central symphysis. This made the lower jaws function as a single connected whole, the mandible. The symphysis was often very thin transversely and long, accounting for a considerable part of the jaw length, up to 60%. If a crest was present on the snout, the symphysis could feature a matching mandible crest, jutting out to below. Toothed species also bore teeth in their dentaries. The mandible opened and closed in a simple vertical or "orthal" up-and-down movement. Vertebral column The vertebral column of pterosaurs numbered between thirty-four and seventy vertebrae. The vertebrae in front of the tail were "procoelous": the cotyle (front of the vertebral body) was concave and into it fitted a convex extension at the rear of the preceding vertebra, the condyle. Advanced pterosaurs are unique in possessing special processes projecting adjacent to their condyle and cotyle, the exapophyses, and the cotyle also may possess a small prong on its midline called a hypapophysis. The necks of pterosaurs were relatively long and straight. In pterodactyloids, the neck is typically longer than the torso. This length is not caused by an increase of the number of vertebrae, which is invariably seven. Some researchers include two transitional "cervicodorsals" which brings the number to nine. Instead, the vertebrae themselves became more elongated, up to eight times longer than wide. Nevertheless, the cervicals were wider than high, implying a better vertical than horizontal neck mobility. Pterodactyloids have lost all neck ribs. Pterosaur necks were probably rather thick and well-muscled, especially vertically. The torso was relatively short and egg-shaped. The vertebrae in the back of pterosaurs originally might have numbered eighteen. With advanced species a growing number of these tended to be incorporated into the sacrum. Such species also often show a fusion of the front dorsal vertebrae into a rigid whole which is called the notarium after a comparable structure in birds. This was an adaptation to withstand the forces caused by flapping the wings. The notarium included three to seven vertebrae, depending on the species involved but also on individual age. These vertebrae could be connected by tendons or a fusion of their neural spines into a "supraneural plate". Their ribs also would be tightly fused into the notarium. In general, the ribs are double headed. The sacrum consisted of three to ten sacral vertebrae. They too, could be connected via a supraneural plate that, however, would not contact the notarium. The tails of pterosaurs were always rather slender. This means that the caudofemoralis retractor muscle which in most basal Archosauria provides the main propulsive force for the hindlimb, was relatively unimportant. The tail vertebrae were amphicoelous, the vertebral bodies on both ends being concave. Early species had long tails, containing up to fifty caudal vertebrae, the middle ones stiffened by elongated articulation processes, the zygapophyses, and chevrons. Such tails acted as rudders, sometimes ending at the rear in a vertical diamond-shaped or oval vane. In pterodactyloids, the tails were much reduced and never stiffened, with some species counting as few as ten vertebrae. Shoulder girdle The shoulder girdle was a strong structure that transferred the forces of flapping flight to the thorax. It was probably covered by thick muscle layers. The upper bone, the shoulder blade, was a straight bar. It was connected to a lower bone, the coracoid that is relatively long in pterosaurs. In advanced species, their combined whole, the scapulocoracoid, was almost vertically oriented. The shoulder blade in that case fitted into a recess in the side of the notarium, while the coracoid likewise connected to the breastbone. This way, both sides together made for a rigid closed loop, able to withstand considerable forces. A peculiarity was that the breastbone connections of the coracoids often were asymmetrical, with one coracoid attached in front of the other. In advanced species the shoulder joint had moved from the shoulder blade to the coracoid. The joint was saddle-shaped and allowed considerable movement to the wing. It faced sideways and somewhat upwards. The breastbone, formed by fused paired sterna, was wide. It had only a shallow keel. Via sternal ribs, it was at its sides attached to the dorsal ribs. At its rear, a row of belly ribs or gastralia was present, covering the entire belly. To the front, a long point, the cristospina, jutted obliquely upwards. The rear edge of the breastbone was the deepest point of the thorax. Clavicles or interclavicles were completely absent. Wings Pterosaur wings were formed by bones and membranes of skin and other tissues. The primary membranes attached to the extremely long fourth finger of each arm and extended along the sides of the body. Where they ended has been very controversial but since the 1990s a dozen specimens with preserved soft tissue have been found that seem to show they attached to the ankles. The exact curvature of the trailing edge, however, is still equivocal. While historically thought of as simple leathery structures composed of skin, research has since shown that the wing membranes of pterosaurs were highly complex dynamic structures suited to an active style of flight. The outer wings (from the tip to the elbow) were strengthened by closely spaced fibers called actinofibrils. The actinofibrils themselves consisted of three distinct layers in the wing, forming a crisscross pattern when superimposed on one another. The function of the actinofibrils is unknown, as is the exact material from which they were made. Depending on their exact composition (keratin, muscle, elastic structures, etc.), they may have been stiffening or strengthening agents in the outer part of the wing. The wing membranes also contained a thin layer of muscle, fibrous tissue, and a unique, complex circulatory system of looping blood vessels. The combination of actinofibrils and muscle layers may have allowed the animal to adjust the wing slackness and camber. As shown by cavities in the wing bones of larger species and soft tissue preserved in at least one specimen, some pterosaurs extended their system of respiratory air sacs into the wing membrane. Parts of the wing The pterosaur wing membrane is divided into three basic units. The first, called the propatagium ("fore membrane"), was the forward-most part of the wing and attached between the wrist and shoulder, creating the "leading edge" during flight. The brachiopatagium ("arm membrane") was the primary component of the wing, stretching from the highly elongated fourth finger of the hand to the hindlimbs. Finally, at least some pterosaur groups had a membrane that stretched between the legs, possibly connecting to or incorporating the tail, called the uropatagium; the extent of this membrane is not certain, as studies on Sordes seem to suggest that it simply connected the legs but did not involve the tail (rendering it a cruropatagium). A common interpretation is that non-pterodactyloid pterosaurs had a broader uro/cruropatagium stretched between their long fifth toes, with pterodactyloids, lacking such toes, only having membranes running along the legs. There has been considerable argument among paleontologists about whether the main wing membranes (brachiopatagia) attached to the hindlimbs, and if so, where. Fossils of the rhamphorhynchoid Sordes, the anurognathid Jeholopterus, and a pterodactyloid from the Santana Formation seem to demonstrate that the wing membrane did attach to the hindlimbs, at least in some species. However, modern bats and flying squirrels show considerable variation in the extent of their wing membranes and it is possible that, like these groups, different species of pterosaur had different wing designs. Indeed, analysis of pterosaur limb proportions shows that there was considerable variation, possibly reflecting a variety of wing-plans. The bony elements of the arm formed a mechanism to support and extend the wing. Near the body, the humerus or upper arm bone is short but powerfully built. It sports a large deltopectoral crest, to which the major flight muscles are attached. Despite the considerable forces exerted on it, the humerus is hollow or pneumatised inside, reinforced by bone struts. The long bones of the lower arm, the ulna and radius, are much longer than the humerus. They were probably incapable of pronation. A bone unique to pterosaurs, known as the pteroid, connected to the wrist and helped to support the forward membrane (the propatagium) between the wrist and shoulder. Evidence of webbing between the three free fingers of the pterosaur forelimb suggests that this forward membrane may have been more extensive than the simple pteroid-to-shoulder connection traditionally depicted in life restorations. The position of the pteroid bone itself has been controversial. Some scientists, notably Matthew Wilkinson, have argued that the pteroid pointed forward, extending the forward membrane and allowing it to function as an adjustable flap. This view was contradicted in a 2007 paper by Chris Bennett, who showed that the pteroid did not articulate as previously thought and could not have pointed forward, but rather was directed inward toward the body as traditionally interpreted. Specimens of Changchengopterus pani and Darwinopterus linglongtaensis show the pteroid in articulation with the proximal syncarpal, suggesting that the pteroid articulated with the 'saddle' of the radiale (proximal syncarpal) and that both the pteroid and preaxial carpal were migrated centralia. The pterosaur wrist consists of two inner (proximal, at the side of the long bones of the arm) and four outer (distal, at the side of the hand) carpals (wrist bones), excluding the pteroid bone, which may itself be a modified distal carpal. The proximal carpals are fused together into a "syncarpal" in mature specimens, while three of the distal carpals fuse to form a distal syncarpal. The remaining distal carpal, referred to here as the medial carpal, but which has also been termed the distal lateral, or pre-axial carpal, articulates on a vertically elongate biconvex facet on the anterior surface of the distal syncarpal. The medial carpal bears a deep concave fovea that opens anteriorly, ventrally and somewhat medially, within which the pteroid articulates, according to Wilkinson. In derived pterodactyloids like pteranodontians and azhdarchoids, metacarpals I-III are small and do not connect to the carpus, instead hanging in contact with the fourth metacarpal. With these derived species, the fourth metacarpal has been enormously elongated, typically equalling or exceeding the length of the long bones of the lower arm. The fifth metacarpal had been lost. In all species, the first to third fingers are much smaller than the fourth, the "wingfinger", and contain two, three and four phalanges respectively. The smaller fingers are clawed, with the ungual size varying among species. In nyctosaurids the forelimb digits besides the wingfinger have been lost altogether. The wingfinger accounts for about half or more of the total wing length. It normally consists of four phalanges. Their relative lengths tend to vary among species, which has often been used to distinguish related forms. The fourth phalanx is usually the shortest. It lacks a claw and has been lost completely by nyctosaurids. It is curved to behind, resulting in a rounded wing tip, which reduces induced drag. The wingfinger is also bent somewhat downwards. When standing, pterosaurs probably rested on their metacarpals, with the outer wing folded to behind. In this position, the "anterior" sides of the metacarpals were rotated to the rear. This would point the smaller fingers obliquely to behind. According to Bennett, this would imply that the wingfinger, able to describe the largest arc of any wing element, up to 175°, was not folded by flexion but by an extreme extension. The wing was automatically folded when the elbow was bowed. A laser-simulated fluorescence scan on Pterodactylus also identified a membranous "fairing" (area conjunctioning the wing with the body at the neck), as opposed to the feathered or fur-composed "fairing" seen in birds and bats respectively. Pelvis The pelvis of pterosaurs was of moderate size compared to the body as a whole. Often the three pelvic bones were fused. The ilium was long and low, its front and rear blades projecting horizontally beyond the edges of the lower pelvic bones. Despite this length, the rod-like form of these processes indicates that the hindlimb muscles attached to them were limited in strength. The, in side view narrow, pubic bone fused with the broad ischium into an ischiopubic blade. Sometimes, the blades of both sides were also fused, closing the pelvis from below and forming the pelvic canal. The hip joint was not perforated and allowed considerable mobility to the leg. It was directed obliquely upwards, preventing a perfectly vertical position of the leg. The front of the pubic bones articulated with a unique structure, the paired prepubic bones. Together these formed a cusp covering the rear belly, between the pelvis and the belly ribs. The vertical mobility of this element suggests a function in breathing, compensating the relative rigidity of the chest cavity. Hindlimbs The hindlimbs of pterosaurs were strongly built, yet relative to their wingspans smaller than those of birds. They were long in comparison to the torso length. The thighbone was rather straight, with the head making only a small angle with the shaft. This implies that the legs were not held vertically below the body but were somewhat sprawling. The shinbone was often fused with the upper ankle bones into a tibiotarsus that was longer than the thighbone. It could attain a vertical position when walking. The calf bone tended to be slender, especially at its lower end that in advanced forms did not reach the ankle, sometimes reducing total length to a third. Typically, it was fused to the shinbone. The ankle was a simple, "mesotarsal", hinge. The, rather long and slender, metatarsus was always splayed to some degree. The foot was plantigrade, meaning that during the walking cycle the sole of the metatarsus was pressed onto the soil. There was a clear difference between early pterosaurs and advanced species regarding the form of the fifth digit. Originally, the fifth metatarsal was robust and not very shortened. It was connected to the ankle in a higher position than the other metatarsals. It bore a long, and often curved, mobile clawless fifth toe consisting of two phalanges. The function of this element has been enigmatic. It used to be thought that the animals slept upside-down like bats, hanging from branches and using the fifth toes as hooks. Another hypothesis held that they stretched the brachiopatagia, but in articulated fossils the fifth digits are always flexed towards the tail. Later it became popular to assume that these toes extended an uropatagium or cruropatagium between them. As the fifth toes were on the outside of the feet, such a configuration would only have been possible if these rotated their fronts outwards in flight. Such a rotation could be caused by an abduction of the thighbone, meaning that the legs would be spread. This would also turn the feet into a vertical position. They then could act as rudders to control yaw. Some specimens show membranes between the toes, allowing them to function as flight control surfaces. The uropatagium or cruropatagium would control pitch. When walking the toes could flex upwards to lift the membrane from the ground. In Pterodactyloidea, the fifth metatarsal was much reduced and the fifth toe, if present, little more than a stub. This suggests that their membranes were split, increasing flight maneuverability. The first to fourth toes were long. They had two, three, four and five phalanges respectively. Often the third toe was longest; sometimes the fourth. Flat joints indicate a limited mobility. These toes were clawed but the claws were smaller than the hand claws. Soft tissues The rare conditions that allowed for the fossilisation of pterosaur remains, sometimes also preserved soft tissues. Modern synchrotron or ultraviolet light photography has revealed many traces not visible to the naked eye. These are often imprecisely called "impressions" but mostly consist of petrifications, natural casts and transformations of the original material. They may include horn crests, beaks or claw sheaths as well as the various flight membranes. Exceptionally, muscles were preserved. Skin patches show small round non-overlapping scales on the soles of the feet, the ankles and the ends of the metatarsals. They covered pads cushioning the impact of walking. Scales are unknown from other parts of the body. Pycnofibers Most or all pterosaurs had hair-like filaments known as pycnofibers on the head and torso. The term "pycnofiber", meaning "dense filament", was coined by palaeontologist Alexander Kellner and colleagues in 2009. Pycnofibers were unique structures similar to, but not homologous (sharing a common origin) with, mammalian hair, an example of convergent evolution. A fuzzy integument was first reported from a specimen of Scaphognathus crassirostris in 1831 by Georg August Goldfuss, but had been widely doubted. Since the 1990s, pterosaur finds and histological and ultraviolet examination of pterosaur specimens have provided incontrovertible proof: pterosaurs had pycnofiber coats. Sordes pilosus (which translates as "hairy demon") and Jeholopterus ninchengensis show pycnofibers on the head and body. The presence of pycnofibers strongly indicates that pterosaurs were endothermic (warm-blooded). They aided thermoregulation, as is common in warm-blooded animals who need insulation to prevent excessive heat-loss. Pycnofibers were flexible, short filaments, about five to seven millimetres long and rather simple in structure with a hollow central canal. Pterosaur pelts might have been comparable in density to many Mesozoic mammals. Relation with feathers Pterosaur filaments could share a common origin with feathers, as speculated in 2002 by Czerkas and Ji. In 2009, Kellner concluded that pycnofibers were structured similarly to theropod proto-feathers. Others were unconvinced, considering the difference with the "quills" found on many of the bird-like maniraptoran specimens too fundamental. A 2018 study of the remains of two small Jurassic-age pterosaurs from Inner Mongolia, China, found that pterosaurs had a wide array of pycnofiber shapes and structures, as opposed to the homogeneous structures that had generally been assumed to cover them. Some of these had frayed ends, very similar in structure to four different feather types known from birds or other dinosaurs but almost never known from pterosaurs prior to the study, suggesting homology. A response to this study was published in 2020, where it was suggested that the structures seen on the anurognathids were actually a result of the decomposition of aktinofibrils: a type of fibre used to strengthen and stiffen the wing. However, in a response to this, the authors of the 2018 paper point to the fact that the presence of the structures extend past the patagium, and the presence of both aktinofibrils and filaments on Jeholopterus ningchengensis and Sordes pilosus. The various forms of filament structure present on the anurognathids in the 2018 study would also require a form of decomposition that would cause the different 'filament' forms seen. They therefore conclude that the most parsimonious interpretation of the structures is that they are filamentous protofeathers. But Liliana D'Alba points out that the description of the preserved integumentary structures on the two anurognathid specimens is still based upon gross morphology. She also points out that Pterorhynchus was described to have feathers to support the claim that feathers had a common origin with Ornithodirans but was argued against by several authors. The only method to assure if it was homologous to feathers is to use a scanning electron microscope. In 2022, a new fossil of Tupandactylus cf. imperator was found to have melanosomes in forms that signal an earlier-than-anticipated development of patterns found in extant feathers. The new specimen suggested that pterosaur integumentary melanosomes exhibited a more complex organization than those previously known from other pterosaurs. This indicates the presence of a unique form of melanosomes within pterosaur integument at the time, distinct from previously known contemporary integumentary structures and more similar to those reported from mammalian hair and avian feathers. The feather fossils obtained from this specimen also suggest the presence of Stage IIIa feathers, a new discovery that indicates more complex feather structures were present in pterosaurs. The study describing this specimen further clarifies the timeline of avian feather evolution and suggests that the feather-specific melanosome signaling found in extant birds are possibly homologous with those found in pterosaurs. History of discovery First finds Pterosaur fossils are very rare, due to their light bone construction. Complete skeletons can generally only be found in geological layers with exceptional preservation conditions, the so-called Lagerstätten. The pieces from one such Lagerstätte, the Late Jurassic Solnhofen Limestone in Bavaria, became much sought after by rich collectors. In 1784, Italian naturalist Cosimo Alessandro Collini was the first scientist to describe a pterosaur fossil. At that time the concepts of evolution and extinction were imperfectly developed. The bizarre build of the pterosaur was shocking, as it could not clearly be assigned to any existing animal group. The discovery of pterosaurs would thus play an important role in the progress of modern paleontology and geology. Scientific opinion at the time was that if such creatures were still alive, only the sea was a credible habitat; Collini suggested it might be a swimming animal that used its long front limbs as paddles. A few scientists continued to support the aquatic interpretation even until 1830, when German zoologist Johann Georg Wagler suggested that Pterodactylus used its wings as flippers and was affiliated with Ichthyosauria and Plesiosauria. In 1800, Johann Hermann first suggested that it represented a flying creature in a letter to Georges Cuvier. Cuvier agreed in 1801, understanding it was an extinct flying reptile. In 1809, he coined the name Ptéro-Dactyle, "wing-finger". This was in 1815 Latinised to Pterodactylus. At first most species were assigned to this genus and ultimately "pterodactyl" was popularly and incorrectly applied to all members of Pterosauria. Today, paleontologists limit the term to the genus Pterodactylus or members of the Pterodactyloidea. In 1812 and 1817, Samuel Thomas von Soemmerring redescribed the original specimen and an additional one. He saw them as affiliated to birds and bats. Although he was mistaken in this, his "bat model" would be influential during the 19th century. In 1843, Edward Newman thought pterosaurs were flying marsupials. Ironically, as the "bat model" depicted pterosaurs as warm-blooded and furred, it would turn out to be more correct in certain aspects than Cuvier's "reptile model" in the long run. In 1834, Johann Jakob Kaup coined the term Pterosauria. Expanding research In 1828, Mary Anning found in England the first pterosaur genus outside Germany, named as Dimorphodon by Richard Owen, also the first non-pterodactyloid pterosaur known. Later in the century, the Early Cretaceous Cambridge Greensand produced thousands of pterosaur fossils, that however, were of poor quality, consisting mostly of strongly eroded fragments. Nevertheless, based on these, numerous genera and species would be named. Many were described by Harry Govier Seeley, at the time the main English expert on the subject, who also wrote the first pterosaur book, Ornithosauria, and in 1901 the first popular book, Dragons of the Air. Seeley thought that pterosaurs were warm-blooded and dynamic creatures, closely related to birds. Earlier, the evolutionist St. George Jackson Mivart had suggested pterosaurs were the direct ancestors of birds. Owen opposed the views of both men, seeing pterosaurs as cold-blooded "true" reptiles. In the US, Othniel Charles Marsh in 1870 discovered Pteranodon in the Niobrara Chalk, then the largest known pterosaur, the first toothless one and the first from America. These layers too rendered thousands of fossils, also including relatively complete skeletons that were three-dimensionally preserved instead of being strongly compressed as with the Solnhofen specimens. This led to a much better understanding of many anatomical details, such as the hollow nature of the bones. Meanwhile, finds from the Solnhofen had continued, accounting for the majority of complete high-quality specimens discovered. They allowed to identify most new basal taxa, such as Rhamphorhynchus, Scaphognathus and Dorygnathus. This material gave birth to a German school of pterosaur research, which saw flying reptiles as the warm-blooded, furry and active Mesozoic counterparts of modern bats and birds. In 1882, Marsh and Karl Alfred Zittel published studies about the wing membranes of specimens of Rhamphorhynchus. German studies continued well into the 1930s, describing new species such as Anurognathus. In 1927, Ferdinand Broili discovered hair follicles in pterosaur skin, and paleoneurologist Tilly Edinger determined that the brains of pterosaurs more resembled those of birds than modern cold-blooded reptiles. In contrast, English and American paleontologists by the middle of the twentieth century largely lost interest in pterosaurs. They saw them as failed evolutionary experiments, cold-blooded and scaly, that hardly could fly, the larger species only able to glide, being forced to climb trees or throw themselves from cliffs to achieve a take-off. In 1914, for the first-time pterosaur aerodynamics were quantitatively analysed, by Ernest Hanbury Hankin and David Meredith Seares Watson, but they interpreted Pteranodon as a pure glider. Little research was done on the group during the 1940s and 1950s. Pterosaur renaissance The situation for dinosaurs was comparable. From the 1960s onwards, a dinosaur renaissance took place, a quick increase in the number of studies and critical ideas, influenced by the discovery of additional fossils of Deinonychus, whose spectacular traits refuted what had become entrenched orthodoxy. In 1970, likewise the description of the furry pterosaur Sordes began what Robert Bakker named a renaissance of pterosaurs. Kevin Padian especially propagated the new views, publishing a series of studies depicting pterosaurs as warm-blooded, active and running animals. This coincided with a revival of the German school through the work of Peter Wellnhofer, who in 1970s laid the foundations of modern pterosaur science. In 1978, he published the first pterosaur textbook, the Handbuch der Paläoherptologie, Teil 19: Pterosauria, and in 1991 the second ever popular science pterosaur book, the Encyclopedia of Pterosaurs. This development accelerated through the exploitation of two new Lagerstätten. During the 1970s, the Early Cretaceous Santana Formation in Brazil began to produce chalk nodules that, though often limited in size and the completeness of the fossils they contained, perfectly preserved three-dimensional pterosaur skeletal parts. German and Dutch institutes bought such nodules from fossil poachers and prepared them in Europe, allowing their scientists to describe many new species and revealing a whole new fauna. Soon, Brazilian researchers, among them Alexander Kellner, intercepted the trade and named even more species. Even more productive was the Early Cretaceous Chinese Jehol Biota of Liaoning that since the 1990s has brought forth hundreds of exquisitely preserved two-dimensional fossils, often showing soft tissue remains. Chinese researchers such as Lü Junchang have again named many new taxa. As discoveries also increased in other parts of the world, a sudden surge in the total of named genera took place. By 2009, when they had increased to about ninety, this growth showed no sign of levelling-off. In 2013, M.P. Witton indicated that the number of discovered pterosaur species had risen to 130. Over ninety percent of known taxa has been named during the "renaissance". Many of these were from groups the existence of which had been unknown. Advances in computing power enabled researchers to determine their complex relationships through the quantitative method of cladistics. New and old fossils yielded much more information when subjected to modern ultraviolet light or roentgen photography, or CAT-scans. Insights from other fields of biology were applied to the data obtained. All this resulted in a substantial progress in pterosaur research, rendering older accounts in popular science books completely outdated. In 2017 a fossil from a 170-million-year-old pterosaur, later named as the species Dearc sgiathanach in 2022, was discovered on the Isle of Skye in Scotland. The National Museum of Scotland claims that it is the largest of its kind ever discovered from the Jurassic period, and it has been described as the world's best-preserved skeleton of a pterosaur. Evolution and extinction Origins Because pterosaur anatomy has been so heavily modified for flight, and immediate transitional fossil predecessors have not so far been described, the ancestry of pterosaurs is not fully understood. The oldest known pterosaurs were already fully adapted to a flying lifestyle. Since Seeley, it was recognised that pterosaurs were likely to have had their origin in the "archosaurs", what today would be called the Archosauromorpha. In the 1980s, early cladistic analyses found that they were Avemetatarsalians (archosaurs closer to dinosaurs than to crocodilians). As this would make them also rather close relatives of the dinosaurs, these results were seen by Kevin Padian as confirming his interpretation of pterosaurs as bipedal warm-blooded animals. Because these early analyses were based on a limited number of taxa and characters, their results were inherently uncertain. Several influential researchers who rejected Padian's conclusions offered alternative hypotheses. David Unwin proposed an ancestry among the basal Archosauromorpha, specifically long-necked forms ("protorosaurs") such as tanystropheids. A placement among basal archosauriforms like Euparkeria was also suggested. Some basal archosauromorphs seem at first glance to be good candidates for close pterosaur relatives due to their long-limbed anatomy; one example is Sharovipteryx, a "protorosaur" with skin membranes on its hindlimbs likely used for gliding. A 1999 study by Michael Benton found that pterosaurs were avemetatarsalians closely related to Scleromochlus, and named the group Ornithodira to encompass pterosaurs and dinosaurs. Two researchers, S. Christopher Bennett in 1996, and paleoartist David Peters in 2000, published analyses finding pterosaurs to be protorosaurs or closely related to them. However, Peters gathered novel anatomical data using an unverified technique called "Digital Graphic Segregation" (DGS), which involves digitally tracing over images of pterosaur fossils using photo editing software. Bennett only recovered pterosaurs as close relatives of the protorosaurs after removing characteristics of the hindlimb from his analysis, to test the possibility of locomotion-based convergent evolution between pterosaurs and dinosaurs. A 2007 reply by Dave Hone and Michael Benton could not reproduce this result, finding pterosaurs to be closely related to dinosaurs even without hindlimb characters. They also criticized David Peters for drawing conclusions without access to the primary evidence, that is, the pterosaur fossils themselves. Hone and Benton concluded that, although more basal pterosauromorphs are needed to clarify their relationships, current evidence indicates that pterosaurs are avemetatarsalians, as either the sister group of Scleromochlus or a branch between the latter and Lagosuchus. A 2011 archosaur-focused phylogenetic analysis by Sterling Nesbitt benefited from far more data and found strong support for pterosaurs being avemetatarsalians, though Scleromochlus was not included due to its poor preservation. A 2016 archosauromorph-focused study by Martin Ezcurra included various proposed pterosaur relatives, yet also found pterosaurs to be closer to dinosaurs and unrelated to more basal taxa. Working from his 1996 analysis, Bennett published a 2020 study on Scleromochlus which argued that both Scleromochlus and pterosaurs were non-archosaur archosauromorphs, albeit not particularly closely related to each other. By contrast, a later 2020 study proposed that lagerpetid archosaurs were the sister clade to pterosauria. This was based on newly described fossil skulls and forelimbs showing various anatomical similarities with pterosaurs and reconstructions of lagerpetid brains and sensory systems based on CT scans also showing neuroanatomical similarities with pterosaurs. The results of the latter study were subsequently supported by an independent analysis of early pterosauromorph interrelationships. A related problem is the origin of pterosaur flight. Like with birds, hypotheses can be ordered into two main varieties: "ground up" or "tree down". Climbing a tree would cause height and gravity to provide both the energy and a strong selection pressure for incipient flight. Rupert Wild in 1983 proposed a hypothetical "propterosaurus": a lizard-like arboreal animal developing a membrane between its limbs, first to safely parachute and then, gradually elongating the fourth finger, to glide. However, subsequent cladistic results did not fit this model well. Neither protorosaurs nor ornithodirans are biologically equivalent to lizards. Furthermore, the transition between gliding and flapping flight is not well-understood. More recent studies on basal pterosaur hindlimb morphology seem to vindicate a connection to Scleromochlus. Like this archosaur, basal pterosaur lineages have plantigrade hindlimbs that show adaptations for saltation. At least one study found that the early Triassic ichnofossil Prorotodactylus is anatomically similar to that of early pterosaurs. Extinction It was once thought that competition with early bird species might have resulted in the extinction of many of the pterosaurs. It was thought that by the end of the Cretaceous, only large species of pterosaurs were present (no longer true; see below). The smaller species were thought to have become extinct, their niche filled by birds. However, pterosaur decline (if actually present) seems unrelated to bird diversity, as ecological overlap between the two groups appears to be minimal. In fact, at least some avian niches were reclaimed by pterosaurs prior to the Cretaceous–Paleogene extinction event. It seems that the K-Pg extinction event at the end of the Cretaceous, which wiped out all non-avian dinosaurs and many other animals, was the direct cause of the extinction of the pterosaurs. In the early 2010s, several new pterosaur taxa were discovered dating to the Campanian/Maastrichtian, such as the ornithocheirids Piksi and "Ornithocheirus", possible pteranodontids and nyctosaurids, several tapejarids and the indeterminate non-azhdarchid Navajodactylus. Small azhdarchoid pterosaurs were also present in the Campanian. This suggests that late Cretaceous pterosaur faunas were far more diverse than previously thought, possibly not even having declined significantly from the early Cretaceous. Small-sized pterosaur species apparently were present in the Csehbánya Formation, indicating a higher diversity of Late Cretaceous pterosaurs than previously accounted for. The recent findings of a small cat-sized adult azhdarchid further indicate that small pterosaurs from the Late Cretaceous might actually have simply been rarely preserved in the fossil record, helped by the fact that there is a strong bias against terrestrial small sized vertebrates such as juvenile dinosaurs, and that their diversity might actually have been much larger than previously thought. At least some non-pterodactyloid pterosaurs survived into the Late Cretaceous, postulating a Lazarus taxa situation for late Cretaceous pterosaur faunas. A 2021 study showcases that niches previously occupied by small pterosaurs were increasingly occupied by the juvenile stages of larger species in the Late Cretaceous. Rather than being outcompeted by birds, pterosaurs essentially specialized a trend already occurring in previous eras of the Mesozoic. Classification and phylogeny In phylogenetic taxonomy, the clade Pterosauria has usually been defined as node-based and anchored to several extensively studied taxa as well as those thought to be primitive. One 2003 study defined Pterosauria as "The most recent common ancestor of the Anurognathidae, Preondactylus and Quetzalcoatlus and all their descendants." However, these types of definition would inevitably leave any related species that are slightly more primitive out of the Pterosauria. To remedy this, a new definition was proposed that would anchor the name not to any particular species but to an anatomical feature, the presence of an enlarged fourth finger that supports a wing membrane. This "apomorophy-based" definition was adopted by the PhyloCode in 2020 as "[T]he clade characterized by the apomorphy fourth manual digit hypertrophied to support a wing membrane, as inherited by Pterodactylus (originally Ornithocephalus) antiquus (Sömmerring 1812)". A broader clade, Pterosauromorpha, has been defined as all ornithodirans more closely related to pterosaurs than to dinosaurs. The internal classification of pterosaurs has historically been difficult, because there were many gaps in the fossil record. Starting from the 21st century, new discoveries are now filling in these gaps and giving a better picture of the evolution of pterosaurs. Traditionally, they were organized into two suborders: the Rhamphorhynchoidea, a "primitive" group of long-tailed pterosaurs, and the Pterodactyloidea, "advanced" pterosaurs with short tails. However, this traditional division has been largely abandoned. Rhamphorhynchoidea is a paraphyletic (unnatural) group, since the pterodactyloids evolved directly from them and not from a common ancestor, so, with the increasing use of cladistics, it has fallen out of favor among most scientists. The precise relationships between pterosaurs is still unsettled. Many studies of pterosaur relationships in the past have included limited data and were highly contradictory. However, newer studies using larger data sets are beginning to make things clearer. The cladogram (family tree) below follows a phylogenetic analysis presented by Longrich, Martill and Andres in 2018, with clade names after Andres et al. (2014). Paleobiology Flight The mechanics of pterosaur flight are not completely understood or modeled at this time. Katsufumi Sato, a Japanese scientist, did calculations using modern birds and concluded that it was impossible for a pterosaur to stay aloft. In the book Posture, Locomotion, and Paleoecology of Pterosaurs it is theorized that they were able to fly due to the oxygen-rich, dense atmosphere of the Late Cretaceous period. However, both Sato and the authors of Posture, Locomotion, and Paleoecology of Pterosaurs based their research on the now-outdated theories of pterosaurs being seabird-like, and the size limit does not apply to terrestrial pterosaurs, such as azhdarchids and tapejarids. Furthermore, Darren Naish concluded that atmospheric differences between the present and the Mesozoic were not needed for the giant size of pterosaurs. Another issue that has been difficult to understand is how they took off. Earlier suggestions were that pterosaurs were largely cold-blooded gliding animals, deriving warmth from the environment like modern lizards, rather than burning calories. In this case, it was unclear how the larger ones of enormous size, with an inefficient cold-blooded metabolism, could manage a bird-like takeoff strategy, using only the hind limbs to generate thrust for getting airborne. Later research shows them instead as being warm-blooded and having powerful flight muscles, and using the flight muscles for walking as quadrupeds. Mark Witton of the University of Portsmouth and Mike Habib of Johns Hopkins University suggested that pterosaurs used a vaulting mechanism to obtain flight. The tremendous power of their winged forelimbs would enable them to take off with ease. Once aloft, pterosaurs could reach speeds of up to and travel thousands of kilometres. In 1985, the Smithsonian Institution commissioned aeronautical engineer Paul MacCready to build a half-scale working model of Quetzalcoatlus northropi. The replica was launched with a ground-based winch. It flew several times in 1986 and was filmed as part of the Smithsonian's IMAX film On the Wing. Large-headed species are thought to have forwardly swept their wings in order to better balance. Air sacs and respiration A 2009 study showed that pterosaurs had a lung-and-air-sac system and a precisely controlled skeletal breathing pump, which supports a flow-through pulmonary ventilation model in pterosaurs, analogous to that of birds. The presence of a subcutaneous air sac system in at least some pterodactyloids would have further reduced the density of the living animal. Like modern crocodilians, pterosaurs appeared to have had a hepatic piston, seeing as their shoulder-pectoral girdles were too inflexible to move the sternum as in birds, and they possessed strong gastralia. Thus, their respiratory system had characteristics comparable to both modern archosaur clades. Nervous system An X-ray study of pterosaur brain cavities revealed that the animals (Rhamphorhynchus muensteri and Anhanguera santanae) had massive flocculi. The flocculus is a brain region that integrates signals from joints, muscles, skin and balance organs. The pterosaurs' flocculi occupied 7.5% of the animals' total brain mass, more than in any other vertebrate. Birds have unusually large flocculi compared with other animals, but these only occupy between 1 and 2% of total brain mass. The flocculus sends out neural signals that produce small, automatic movements in the eye muscles. These keep the image on an animal's retina steady. Pterosaurs may have had such a large flocculus because of their large wing size, which would mean that there was a great deal more sensory information to process. The low relative mass of the flocculi in birds is also a result of birds having a much larger brain overall; though this has been considered an indication that pterosaurs lived in a structurally simpler environment or had less complex behaviour compared to birds, recent studies of crocodilians and other reptiles show that it is common for sauropsids to achieve high intelligence levels with small brains. Studies on the endocast of Allkaruen show that brain evolution in pterodactyloids was a modular process. Terrestrial locomotion Pterosaurs' hip sockets are oriented facing slightly upwards, and the head of the femur (thigh bone) is only moderately inward facing, suggesting that pterosaurs had an erect stance. It would have been possible to lift the thigh into a horizontal position during flight, as gliding lizards do. There was considerable debate whether pterosaurs ambulated as quadrupeds or as bipeds. In the 1980s, paleontologist Kevin Padian suggested that smaller pterosaurs with longer hindlimbs, such as Dimorphodon, might have walked or even run bipedally, in addition to flying, like road runners. However, a large number of pterosaur trackways were later found with a distinctive four-toed hind foot and three-toed front foot; these are the unmistakable prints of pterosaurs walking on all fours. Fossil footprints show that pterosaurs stood with the entire foot in contact with the ground (plantigrade), in a manner similar to many mammals like humans and bears. Footprints from azhdarchids and several unidentified species show that pterosaurs walked with an erect posture with their four limbs held almost vertically beneath the body, an energy-efficient stance used by most modern birds and mammals, rather than the sprawled limbs of modern reptiles. Indeed, erect-limbs may be omnipresent in pterosaurs. Though traditionally depicted as ungainly and awkward when on the ground, the anatomy of some pterosaurs (particularly pterodactyloids) suggests that they were competent walkers and runners. Early pterosaurs have long been considered particularly cumbersome locomotors due to the presence of large cruropatagia, but they too appear to have been generally efficient on the ground. The forelimb bones of azhdarchids and ornithocheirids were unusually long compared to other pterosaurs, and, in azhdarchids, the bones of the arm and hand (metacarpals) were particularly elongated. Furthermore, as a whole, azhdarchid front limbs were proportioned similarly to fast-running ungulate mammals. Their hind limbs, on the other hand, were not built for speed, but they were long compared with most pterosaurs, and allowed for a long stride length. While azhdarchid pterosaurs probably could not run, they would have been relatively fast and energy efficient. The relative size of the hands and feet in pterosaurs (by comparison with modern animals such as birds) may indicate the type of lifestyle pterosaurs led on the ground. Azhdarchid pterosaurs had relatively small feet compared to their body size and leg length, with foot length only about 25–30% the length of the lower leg. This suggests that azhdarchids were better adapted to walking on dry, relatively solid ground. Pteranodon had slightly larger feet (47% the length of the tibia), while filter-feeding pterosaurs like the ctenochasmatoids had very large feet (69% of tibial length in Pterodactylus, 84% in Pterodaustro), adapted to walking in soft muddy soil, similar to modern wading birds. Though clearly forelimb-based launchers, basal pterosaurs have hindlimbs well adapted for hopping, suggesting a connection with archosaurs such as Scleromochlus. Swimming Tracks made by ctenochasmatoids indicate that these pterosaurs swam using their hindlimbs. In general, these have large hindfeet and long torsos, indicating that they were probably more adapted for swimming than other pterosaurs. Pteranodontians conversely have several speciations in their humeri interpreted to have been suggestive of a water-based version of the typical quadrupedal launch, and several like boreopterids must have foraged while swimming, as they seem incapable of frigatebird-like aerial hawking. These adaptations are also seen in terrestrial pterosaurs like azhdarchids, which presumably still needed to launch from water in case they found themselves in it. The nyctosaurid Alcione may display adaptations for wing-propelled diving like modern gannets and tropicbirds. Diet and feeding habits Traditionally, almost all pterosaurs were seen as surface-feeding piscivores or fish-eaters, a view that still dominates popular science. Today, many pterosaurs groups are thought to have been terrestrial carnivores, omnivores or insectivores. Early-on it was recognised that the small Anurognathidae were nocturnal, aerial insectivores. With highly flexible joints on the wing finger, a broad, triangular wing shape, large eyes and short tail, these pterosaurs were likely analogous to nightjars or extant insectivorous bats, being capable of high manoeuvrability at relatively low speeds. Interpretations of the habits of basal groups have changed profoundly. Dimorphodon, envisioned as a puffin analogue in the past, is indicated by its jaw structure, gait, and poor flight capabilities, as a terrestrial/semiarboreal predator of small mammals, squamates, and large insects. Its robust dentition caused Campylognathoides to be seen as a generalist or a terrestrial predator of small vertebrates, but the highly robust humerus and high-aspect wing morphology, suggest it may have been capable of grabbing prey on the wing; a later study indicates it was teuthophagous based on squid findings within its gut. The small insectivorous Carniadactylus and the larger Eudimorphodon were highly aerial animals and fast, agile flyers with long robust wings. Eudimorphodon has been found with fish remains in its stomach, but its dentition suggests an opportunistic diet. Slender-winged Austriadactylus and Caviramus were likely terrestrial/semiarboreal generalists. Caviramus likely had a strong bite force, indicating an adaptation towards hard food items that might have been chewed in view of the tooth wear. Some Rhamphorhynchidae, such as Rhamphorhynchus itself or Dorygnathus, were fish-eaters with long, slender wings, needle-like dentition and long, thin jaws. Sericipterus, Scaphognathus and Harpactognathus had more robust jaws and teeth (which were ziphodont, dagger-shaped, in Sericipterus), and shorter, broader wings. These were either terrestrial/aerial predators of vertebrates or corvid-like generalists. Wukongopteridae like Darwinopterus were first considered aerial predators. Lacking a robust jaw structure or powerful flying muscles, they are now seen as arboreal or semiterrestrial insectivores. Darwinopterus robustidens, in particular, seems to have been a beetle specialist. Among pterodactyloids, a greater variation in diet is present. Pteranodontia contained many piscivorous taxa, such as the Ornithocheirae, Boreopteridae, Pteranodontidae and Nyctosauridae. Niche partitioning caused ornithocheirans and the later nyctosaurids to be aerial dip-feeders like today's frigatebirds (with the exception of the plunge-diving adapted Alcione elainus), while boreopterids were freshwater diving animals similar to cormorants, and pteranodonts pelagic plunge-divers akin to boobies and gannets. An analysis of Lonchodraco found clusters of foramina at the tip of its beak; birds with similarly numerous foramina have sensitive beaks used to feel for food, so Lonchodraco may have used its beak to feel for fish or invertebrates in shallow water. The istiodactylids were likely primarily scavengers. Archaeopterodactyloidea obtained food in coastal or freshwater habitats. Germanodactylus and Pterodactylus were piscivores, while the Ctenochasmatidae were suspension feeders, using their numerous fine teeth to filter small organisms from shallow water. Pterodaustro was adapted for flamingo-like filter-feeding. In contrast, Azhdarchoidea mostly were terrestrial pterosaurs. Tapejaridae were arboreal omnivores, supplementing seeds and fruits with small insects and vertebrates. Dsungaripteridae were specialist molluscivores, using their powerful jaws to crush the shells of molluscs and crustaceans. Thalassodromidae were likely terrestrial carnivores. Thalassodromeus itself was named after a fishing method known as "skim-feeding", later understood to be biomechanically impossible. Perhaps it pursued relatively large prey, in view of its reinforced jaw joints and relatively high bite force. Azhdarchidae are now understood to be terrestrial predators akin to ground hornbills or some storks, eating any prey item they could swallow whole. Hatzegopteryx was a robustly built predator of relatively large prey, including medium-sized dinosaurs. Alanqa may have been a specialist molluscivore. A 2021 study reconstructed the adductor musculature of skulls from pterodactyloids, estimating the bite force and potential dietary habits of nine selected species. The study corroborated the view of pteranodontids, nyctosaurids and anhanuerids as piscivores based on them being relatively weak but fast biters, and suggest that Tropeognathus mesembrinus was specialised in consuming relatively large prey compared to Anhanguera. Dsungaripterus was corroborated as a durophage, with Thalassodromeus proposed to share this feeding habit based on high estimated bite force quotients (BFQ) and absolute bite force values. Tapejara wellnhoferi was corroborated as a specialised consumer of hard plant material with a relatively high BFQ and high mechanical advantage, and Caupedactylus ybaka and Tupuxuara leonardii were proposed to be ground-feeding generalists with intermediate bite force values and less specialised jaws. Natural predators Pterosaurs are known to have been eaten by theropods. In the 1 July 2004 edition of Nature, paleontologist Éric Buffetaut discusses an Early Cretaceous fossil of three cervical vertebrae of a pterosaur with the broken tooth of a spinosaur, most likely Irritator, embedded in it. The vertebrae are known not to have been eaten and exposed to digestion, as the joints are still articulated. Fossils of Pteranodon have been found with tooth marks from sharks such as Squalicorax, and a fossil with tooth marks from the Toolebuc formation has been interpreted as being attacked or scavenged by an ichthyosaur (most likely Platypterygius). Reproduction and life history While very little is known about pterosaur reproduction, it is believed that, similar to all dinosaurs, all pterosaurs reproduced by laying eggs, though such findings are very rare. The first known pterosaur eggs were found in the quarries of Liaoning, the same place that yielded feathered dinosaurs, and in Loma del Pterodaustro (Lagarcito Formation, Argentina). The eggs from Liaoning were squashed flat with no signs of cracking, so evidently the eggs had leathery shells, as in modern lizards. The egg from the Lagarcito Formation was laid by a Pterodaustro, a pterosaur known by abundant material. This was supported by the description of an additional pterosaur egg belonging to the genus Darwinopterus, described in 2011, which also had a leathery shell and, also like modern reptiles but unlike birds, was fairly small compared to the size of the mother. In 2014 five unflattened eggs from the species Hamipterus tianshanensis were found in an Early Cretaceous deposit in northwest China. Examination of the shells by scanning electron microscopy showed the presence of a thin calcareous eggshell layer with a membrane underneath. A study of pterosaur eggshell structure and chemistry published in 2007 indicated that it is likely pterosaurs buried their eggs, like modern crocodiles and turtles. Egg-burying would have been beneficial to the early evolution of pterosaurs, as it allows for more weight-reducing adaptations, but this method of reproduction would also have put limits on the variety of environments pterosaurs could live in and may have disadvantaged them when they began to face ecological competition from birds. A Darwinopterus specimen showcases that at least some pterosaurs had a pair of functional ovaries, as opposed to the single functional ovary in birds, dismissing the reduction of functional ovaries as a requirement for powered flight. Wing membranes preserved in pterosaur embryos are well developed, suggesting that pterosaurs were ready to fly soon after birth. However, tomography scans of fossilised Hamipterus eggs suggests that the young pterosaurs had well-developed thigh bones for walking, but weak chests for flight. It is unknown if this holds true for other pterosaurs. Fossils of pterosaurs only a few days to a week old (called "flaplings") have been found, representing several pterosaur families, including pterodactylids, rhamphorhinchids, ctenochasmatids and azhdarchids. All preserve bones that show a relatively high degree of hardening (ossification) for their age, and wing proportions similar to adults. In fact, many pterosaur flaplings have been considered adults and placed in separate species in the past. Additionally, flaplings are normally found in the same sediments as adults and juveniles of the same species, such as the Pterodactylus and Rhamphorhynchus flaplings found in the Solnhofen limestone of Germany, and Pterodaustro flaplings from Argentina. All are found in deep aquatic environment far from shore. For the majority of pterosaur species, it is not known whether they practiced any form of parental care, but their ability to fly as soon as they emerged from the egg and the numerous flaplings found in environments far from nests and alongside adults has led most researchers, including Christopher Bennett and David Unwin, to conclude that the young were dependent on their parents for a relatively short period of time, during a period of rapid growth while the wings grew long enough to fly, and then left the nest to fend for themselves, possibly within days of hatching. Alternatively, they may have used stored yolk products for nourishment during their first few days of life, as in modern reptiles, rather than depend on parents for food. Fossilised Hamipterus nests were shown preserving many male and female pterosaurs together with their eggs in a manner to a similar to that of modern seabird colonies. Due to how underdeveloped the chests of the hatchlings were for flying, it was suggested that Hamipterus may have practiced some form of parental care. However, this study has since been criticised. Most evidence currently leans towards pterosaur hatchlings being superprecocial, similar to that of megapode birds, which fly after hatching without the need of parental care. A further study compares evidence for superprecociality and "late term flight" and overwhelmingly suggests that most if not all pterosaurs were capable of flight soon after hatching. A later study suggested that while smaller-bodied pterosaurs were most likely superprecocial or precocial, owing to the consistent or decreasing wing aspect ratio during growth, certain large-bodied pterosaurs, such as Pteranodon showed possible evidence of their young being altricial, due to the fast rate the limb bones closest to the body grew compared to any other element of their skeleton after hatching. Other factors mentioned were the limits of soft shelled eggs and the size of the pelvic opening of large female pterosaurs. Growth rates of pterosaurs once they hatched varied across different groups. In more primitive, long-tailed pterosaurs ("rhamphorhynchoids"), such as Rhamphorhynchus, the average growth rate during the first year of life was 130% to 173%, slightly faster than the growth rate of alligators. Growth in these species slowed after sexual maturity, and it would have taken more than three years for Rhamphorhynchus to attain maximum size. In contrast, the more advanced, large pterodactyloid pterosaurs, such as Pteranodon, grew to adult size within the first year of life. Additionally, pterodactyloids had determinate growth, meaning that the animals reached a fixed maximum adult size and stopped growing. A 2021 study indicates that pterosaur juveniles of larger species increasingly took the roles previously occupied by adult small pterosaurs. Daily activity patterns Comparisons between the scleral rings of pterosaurs and modern birds and reptiles have been used to infer daily activity patterns of pterosaurs. The pterosaur genera Pterodactylus, Scaphognathus, and Tupuxuara have been inferred to be diurnal, Ctenochasma, Pterodaustro, and Rhamphorhynchus have been inferred to be nocturnal, and Tapejara has been inferred to be cathemeral, being active throughout the day for short intervals. As a result, the possibly fish-eating Ctenochasma and Rhamphorhynchus may have had similar activity patterns to modern nocturnal seabirds, and the filter-feeding Pterodaustro may have had similar activity patterns to modern anseriform birds that feed at night. The differences between activity patterns of the Solnhofen pterosaurs Ctenochasma, Rhamphorhynchus, Scaphognathus, and Pterodactylus may also indicate niche partitioning between these genera. Cultural significance Pterosaurs have been a staple of popular culture for as long as their cousins the dinosaurs, though they are usually not featured as prominently in films, literature or other art. While the depiction of dinosaurs in popular media has changed radically in response to advances in paleontology, a mainly outdated picture of pterosaurs has persisted since the mid-20th century. The vague generic term "pterodactyl" is often used for these creatures. The animals depicted in fiction and pop culture frequently represent either the Pteranodon or (non-pterodactyloid) Rhamphorhynchus, or a fictionalized hybrid of the two. Many children's toys and cartoons feature "pterodactyls" with Pteranodon-like crests and long, Rhamphorhynchus-like tails and teeth, a combination that never existed in nature. However, at least one pterosaur did have both the Pteranodon-like crest and teeth: Ludodactylus, whose name means "toy finger" for its resemblance to old, inaccurate children's toys. Pterosaurs have sometimes been incorrectly identified as (the ancestors of) birds, though birds are theropod dinosaurs and not descendants of pterosaurs. Pterosaurs were used in fiction in Sir Arthur Conan Doyle's 1912 novel The Lost World and its 1925 film adaptation. They appeared in a number of films and television programs since, including the 1933 film King Kong, and 1966's One Million Years B.C. In the latter, animator Ray Harryhausen had to add inaccurate bat-like wing fingers to his stop motion models in order to keep the membranes from falling apart, though this particular error was common in art even before the film was made. Rodan, a fictional giant monster (or kaiju) which first appeared in the 1956 film Rodan, is portrayed as an enormous irradiated species of Pteranodon. Rodan has appeared in multiple Japanese Godzilla films released during the 1960s, 1970s, 1990s, and 2000s, and also appeared in the 2019 American-produced film Godzilla: King of the Monsters. After the 1960s, pterosaurs remained mostly absent from notable American film appearances until 2001's Jurassic Park III. Paleontologist Dave Hone noted that the pterosaurs in this film had not been significantly updated to reflect modern research. Errors persisting were teeth while toothless Pteranodon was intended to be depicted, nesting behavior that was known to be inaccurate by 2001, and leathery wings, rather than the taut membranes of muscle fiber required for pterosaur flight. Petrie from The Land Before Time (1988), is a notable example from an animated film. In most media appearances, pterosaurs are depicted as piscivores, not reflecting their full dietary variation. They are also often shown as aerial predators similar to birds of prey, grasping human victims with talons on their feet. However, only the small anurognathid Vesperopterylus and small wukongopterid Kunpengopterus are known to possess prehensile feet and hands respectively; all other known pterosaurs have flat, plantigrade feet with no opposable toes, and the feet are generally proportionally small, at least in the case of the Pteranodontia.	Biology and health sciences	Dinosaurs and prehistoric reptiles	null
24825	https://en.wikipedia.org/wiki/Pteranodon	Pteranodon	Pteranodon (; from and ) is a genus of pterosaur that included some of the largest known flying reptiles, with P. longiceps having a wingspan of over . They lived during the late Cretaceous geological period of North America in present-day Kansas, Nebraska, Wyoming, South Dakota and Alabama. More fossil specimens of Pteranodon have been found than any other pterosaur, with about 1,200 specimens known to science, many of them well preserved with nearly complete skulls and articulated skeletons. It was an important part of the animal community in the Western Interior Seaway. When the first fossils of Pteranodon were found, they were assigned to toothed pterosaur genera, Ornithocheirus and Pterodactylus. In 1876, Othniel Charles Marsh recognised it as a genus of its own, making particular note of its complete lack of teeth, which at the time was unique among pterosaurs. Over the decades, multiple species would be assigned to Pteranodon, though today, only two are recognised: P. longiceps, the type species, and P. sternbergi. A third species, P. maiseyi, may also exist. Some have suggested that the latter two as genus of their own, Geosternbergia, though this is the subject of some debate. Another genus split from Pteranodon, Dawndraco, may be synonymous with Geosternbergia if that genus is valid, or with Pteranodon if it is not. Pteranodon is part of the family Pteranodontidae, part of the clade Pteranodontia, which also includes nyctosaurids. Pteranodontians form a larger clade, Pteranodontoidea, alongside ornithocheiromorphs, and that clade falls under the suborder Pterodactyloidea. While not dinosaurs, pterosaurs such as Pteranodon form a clade closely related to dinosaurs as both fall within the clade Avemetatarsalia. Male and female Pteranodon differed in size and crest shape. Males attained wingspans of ; females were smaller, averaging . The crests of males were far larger than those of females. In P. longiceps, they were long and backswept, whereas in P. sternbergi, they were tall and upright. Females also had wider pelvises than males. Discovery and history First fossils Pteranodon was the first pterosaur found outside of Europe. Its fossils first were found by Othniel Charles Marsh in 1871, in the Late Cretaceous Smoky Hill Chalk deposits of western Kansas. These chalk beds were deposited at the bottom of what was once the Western Interior Seaway, a large shallow sea over what now is the midsection of the North American continent. These first specimens, YPM 1160 and YPM 1161, consisted of partial wing bones, as well as a tooth from the prehistoric fish Xiphactinus, which Marsh mistakenly believed to belong to this new pterosaur (all known pterosaurs up to that point had teeth). In 1871, Marsh named the find Pterodactylus oweni, assigning it to the well-known (but much smaller) European genus Pterodactylus. Marsh also collected more wing bones of the large pterosaur in 1871. Realizing that the name he had chosen had already been used for Harry Seeley's European pterosaur species Pterodactylus oweni in 1864, Marsh renamed his giant North American pterosaur Pterodactylus occidentalis, meaning "Western wing finger," in his 1872 description of the new specimen. He named two additional species, based on size differences: Pterodactylus ingens (the largest specimen so far), and Pterodactylus velox (the smallest). Meanwhile, Marsh's rival Edward Drinker Cope had unearthed several specimens of the large North American pterosaur. Based on these specimens, Cope named two new species, Ornithochirus umbrosus and Ornithochirus harpyia, in an attempt to assign them to the large European genus Ornithocheirus, though he misspelled the name (forgetting the 'e'). Cope's paper naming his species was published in 1872, just five days after Marsh's paper. This resulted in a dispute, fought in the published literature, over whose names had priority in what obviously were the same species. Cope conceded in 1875 that Marsh's names did have priority over his, but maintained that Pterodactylus umbrosus was a distinct species (but not genus) from any that Marsh had named previously. Re-evaluation by later scientists has supported Marsh's case, refuting Cope's assertion that P. umbrosus represented a larger, distinct species. A toothless pterosaur While the first Pteranodon wing bones were collected by Marsh and Cope in the early 1870s, the first Pteranodon skull was found on May 2, 1876, along the Smoky Hill River in Wallace County (now Logan County), Kansas, USA, by Samuel Wendell Williston, a fossil collector working for Marsh. A second, smaller skull soon was discovered as well. These skulls showed that the North American pterosaurs were different from any European species, in that they lacked teeth and had bony crests on their skulls. Marsh recognized this major difference, describing the specimens as "distinguished from all previously known genera of the order Pterosauria by the entire absence of teeth." Marsh recognized that this characteristic warranted a new genus, and he coined the name Pteranodon ("wing without tooth") in 1876. Marsh reclassified all the previously named North American species from Pterodactylus to Pteranodon. He considered the smaller skull to belong to Pteranodon occidentalis, based on its size. Marsh classified the larger skull, YPM 1117, in the new species Pteranodon longiceps, which he thought to be a medium-sized species in between the small P. occidentalis and the large P. ingens. Marsh also named several additional species: Pteranodon comptus and Pteranodon nanus were named for fragmentary skeletons of small individuals, while Pteranodon gracilis was based on a wing bone that he mistook for a pelvic bone. He soon realized his mistake, and re-classified that specimen again into a separate genus, which he named Nyctosaurus. P. nanus was also later recognized as a Nyctosaurus specimen. In 1892, Samuel Williston examined the question of Pteranodon classification. He noticed that, in 1871, Seeley had mentioned the existence of a partial set of toothless pterosaur jaws from the Cambridge Greensand of England, which he named Ornithostoma. Because the primary characteristic Marsh had used to separate Pteranodon from other pterosaurs was its lack of teeth, Williston concluded that "Ornithostoma" must be considered the senior synonym of Pteranodon. However, in 1901, Pleininger pointed out that "Ornithostoma" had never been scientifically described or even assigned a species name until Williston's work, and therefore had been a nomen nudum and could not beat out Pteranodon for naming priority. Williston accepted this conclusion and went back to calling the genus Pteranodon. However, both Williston and Pleininger were incorrect, because unnoticed by both of them was the fact that, in 1891, Seeley himself had finally described and properly named Ornithostoma, assigning it to the species O. sedgwicki. In the 2010s, more research on the identity of Ornithostoma showed that it was probably not Pteranodon or even a close relative, but may in fact have been an azhdarchoid, a different type of toothless pterosaur. Revising species Williston was also the first scientist to critically evaluate all of the Pteranodon species classified by Cope and Marsh. He agreed with most of Marsh's classification, with a few exceptions. First, he did not believe that P. ingens and P. umbrosus could be considered synonyms, which even Cope had come to believe. He considered both P. velox and P. longiceps to be dubious; the first was based on non-diagnostic fragments, and the second, though known from a complete skull, probably belonged to one of the other, previously-named species. In 1903, Williston revisited the question of Pteranodon classification, and revised his earlier conclusion that there were seven species down to just three. He considered both P. comptus and P. nanus to be specimens of Nyctosaurus, and divided the others into small (P. velox), medium (P. occidentalis), and large species (P. ingens), based primarily on the shape of their upper arm bones. He thought P. longiceps, the only one known from a skull, could be a synonym of either P. velox or P. occidentalis, based on its size. In 1910, Eaton became the first scientist to publish a more detailed description of the entire Pteranodon skeleton, as it was known at the time. He used his findings to revise the classification of the genus once again based on a better understanding of the differences in pteranodont anatomy. Eaton conducted experiments using clay models of bones to help determine the effects of crushing and flattening on the shapes of the arm bones Williston had used in his own classification. Eaton found that most of the differences in bone shapes could be easily explained by the pressures of fossilization, and concluded that no Pteranodon skeletons had any significant differences from each other besides their size. Therefore, Eaton was left to decide his classification scheme based on differences in the skulls alone, which he assigned to species just as Marsh did, by their size. In the end, Eaton recognized only three valid species: P. occidentalis, P. ingens, and P. longiceps. The discovery of specimens with upright crests, classified by Harksen in 1966 as the new species Pteranodon sternbergi, complicated the situation even further. prompting another revision of the genus by Halsey W. Miller in 1972. Because it was impossible to determine crest shape for all of the species based on headless skeletons, Miller concluded that all Pteranodon species except the two based on skulls (P. longiceps and P. sternbergi) must be considered nomena dubia and abandoned. The skull Eaton thought belonged to P. ingens was placed in the new species Pteranodon marshi, and the skull Eaton assigned to P. occidentalis was re-named Pteranodon eatoni. Miller also recognized another species based on a skull with a crest similar to that of P. sternbergi; Miller named this Pteranodon walkeri. To help bring order to this tangle of names, Miller created three subgenera. P. marshi and P. longiceps were placed in the subgenus Longicepia, though this was later changed to simply Pteranodon due to the rules of priority. P. sternbergi and P. walkeri, the upright-crested species, were given the subgenus Sternbergia, which was later changed to Geosternbergia because Sternbergia was preoccupied. Finally, Miller named the subgenus Occidentalia for P. eatoni, the skull formerly associated with P. occidentalis. Miller further expanded the concept of Pteranodon to include Nyctosaurus as a fourth subgenus. Miller considered these to be an evolutionary progression, with the primitive Nyctosaurus, at the time thought to be crestless, giving rise to small-crested Occidentalia, which in turn gave rise to long-crested Pteranodon, finally leading to tall-crested Geosternbergia. However, Miller made several mistakes in his study concerning which specimens Marsh had assigned to which species, and most scientists disregarded his work on the subject in their later research. In 1984, Robert Milton Schoch published another revision that essentially returned to Marsh's original classification scheme, most notably sinking P. longiceps as a synonym of P. ingens. Recognizing variation In the late 1980s and early 1990s, S. Christopher Bennett published several major papers reviewing the anatomy, taxonomy and life history of Pteranodon. In 1992, he published a paper discussing sexual dimorphism and its role in individual variation among Pteranodon fossils, a follow-up of a 1987 paper he authored on the same subject. In the 1992 paper, he referred only to two species, P. longiceps and P. sternbergi. Two years later, he published a paper fully revising its taxonomy, wherein he concluded that only P. longiceps and P. sternbergi were valid species. P. marshi and P. walkeri were regarded as junior synonyms of P. longiceps, and P. eatoni as a junior synonym of P. stenbergi. The remainder were either rendered nomina dubia or placed in Nyctosaurus. Description Body size and sexual dimorphism Adult male Pteranodon were among the largest pterosaurs, and were the largest flying animals known until the late 20th century, when the giant azhdarchid pterosaurs were discovered. The wingspan of an average adult male Pteranodon was . Adult females were much smaller, averaging in wingspan. A large specimen of Pteranodon longiceps, USNM 50130, is estimated to have a wingspan of , body length of and body mass of . Even larger specimens had wingspans of . Size aside, females were distinguished by their short, rounded head crests and wide pelvic canals, whereas had narrow hips and very large head crests, likely serving a display function. Methods used to estimate the mass of large male Pteranodon specimens (those with wingspans of about 7 meters) have been notoriously unreliable, producing a wide range of estimates. In a review of pterosaur size estimates published in 2010, Mark Witton and Michael Habib argued that the largest estimate of is much too high and an upper limit of is more realistic. Witton and Habib considered the methods used by researchers who obtained smaller mass estimates equally flawed. Most have been produced by scaling modern animals such as bats and birds up to Pteranodon size, despite the fact that pterosaurs have vastly different body proportions and soft tissue anatomy from any living animal. Skull and beak Unlike earlier pterosaurs, such as Rhamphorhynchus and Pterodactylus, Pteranodon had toothless beaks, similar to those of birds. Pteranodon beaks were made of solid, bony margins that projected from the base of the jaws. The beaks were long, slender, and ended in thin, sharp points. The upper jaw, which was longer than the lower jaw, was curved upward; while this normally has been attributed only to the upward-curving beak, one specimen (UALVP 24238) has a curvature corresponding with the beak widening towards the tip. While the tip of the beak is not known in this specimen, the level of curvature suggests it would have been extremely long. The unique form of the beak in this specimen led Alexander Kellner to assign it to a distinct genus, Dawndraco, in 2010. The most distinctive characteristic of Pteranodon is its cranial crest. These crests consisted of skull bones (frontals) projecting upward and backward from the skull. The size and shape of these crests varied due to a number of factors, including age, sex, and species. Male Pteranodon sternbergi, the older species of the two described to date, had a more vertical crest with a broad forward projection, while their descendants, Pteranodon longiceps, evolved a narrower, more backward-projecting crest. Females of both species were smaller and bore small, rounded crests. The crests were probably mainly display structures, though they may have had other functions as well. Postcranial skeleton The neural spines of Pteranodon's vertebrae were narrow. Like many pterosaurs and birds, it possessed a notarium, a fused mass comprising the first six dorsal vertebrae. Similarly, the first few ribs were fused. The pelvic bones were fused to the synsacrum, a mass of vertebrae that included at least two dorsal vertebrae, the sacral vertebrae, and the first caudal vertebra. The sacrals were strengthened by bony ligaments. Beyond the synsacrum, the tail was relatively short, and the last few vertebrae were fused into a bony rod. The entire length of the tail was about 3.5% as long as the wingspan, or up to in the largest males. Pteranodon's scapulae were oriented in such a way that each one braces the other, due to their fusion with the coracoids, providing increased integrity during flight. The humeri were extremely robust, with large, curved deltopectoral crests. The radius and ulna were similarly robust. The first three metacarpals were very slender, and their respective digits sported short, curved unguals (claws). Pteranodon's hind feet had four metatarsals, which were tipped with less curved claws. Paleobiology Flight The wing shape of Pteranodon suggests that it would have flown rather like a modern-day albatross. This is based on the fact that Pteranodon had a high aspect ratio (wingspan to chord length) similar to that of the albatross — 9:1 for Pteranodon, compared to 8:1 for an albatross. Albatrosses spend long stretches of time at sea fishing, and use a flight pattern called "dynamic soaring" which exploits the vertical gradient of wind speed near the ocean surface to travel long distances without flapping, and without the aid of thermals (which do not occur over the open ocean the same way they do over land). While most of a Pteranodon flight would have depended on soaring, like long-winged seabirds, it probably required an occasional active, rapid burst of flapping, and studies of Pteranodon wing loading (the strength of the wings vs. the weight of the body) indicate that they were capable of substantial flapping flight, contrary to some earlier suggestions that they were so big they could only glide. However, a more recent study suggests that it relied on thermal soaring, unlike modern seabirds but much like modern continental flyers and the extinct Pelagornis. Like other pterosaurs, Pteranodon probably took off from a standing, quadrupedal position. Using their long forelimbs for leverage, they would have vaulted themselves into the air in a rapid leap. Almost all of the energy would have been generated by the forelimbs. The upstroke of the wings would have occurred when the animal cleared the ground followed by a rapid down-stroke to generate additional lift and complete the launch into the air. It is possible that Pteranodon could have achieved this from the water, as well as on land, which has been speculated for various other such as the distantly related Anhanguera. Locomotion Historically, terrestrial locomotion in Pteranodon, as in pterosaurs overall, has been the subject of debate, chiefly the matter of whether or not they were bipedal or quadrupedal. The earliest model of Pteranodon locomotion, put forward by Cherrie D. Bramwell and G. R. Whitfield, suggested that they were utterly incapable of walking or standing. Instead, they suggested that it moved on land by pushing itself around, and that it took off by perching on cliffsides and allowing the wind to take it. Subsequent works largely revolved around more conventional methods of locomotion, such as bipedalism and various kinds of quadrupedalism. In 2004, Sankar Chatterjee and R. J. Templin proposed a dual system, wherein pterosaurs walked quadrupedally most of the time, but opted for a bipedal takeoff. The latter, however, is unlikely. Trackways suggest that pterosaurs like Pteranodon were quadrupedal. Diet The diet of Pteranodon is known to have included fish; fossilized fish bones have been found in the stomach area of one Pteranodon, and a fossilized fish bolus has been found between the jaws of another Pteranodon, specimen AMNH 5098. Numerous other specimens also preserve fragments of fish scales and vertebrae near the torso, indicating that fish made up a majority of the diet of Pteranodon (though they may also have taken invertebrates). Traditionally, most researchers have suggested that Pteranodon would have taken fish by dipping their beaks into the water while in low, soaring flight. However, this was probably based on the assumption that the animals could not take off from the water surface. It is more likely that Pteranodon could take off from the water, and would have dipped for fish while swimming rather than while flying. Even a small, female Pteranodon could have reached a depth of at least with its long bill and neck while floating on the surface, and they may have reached even greater depths by plunge-diving into the water from the air like some modern long-winged seabirds. In 1994, Bennett noted that the head, neck, and shoulders of Pteranodon were as heavily built as diving birds, and suggested that they could dive by folding back their wings like the modern gannet. Crest function Pteranodon was notable for its skull crest, though the function of this crest has been a subject of debate. Most explanations have focused on the blade-like, backward pointed crest of male P. longiceps, however, and ignored the wide range of variation across age and sex. The fact that the crests vary so much rules out most practical functions other than for use in mating displays. Therefore, display was probably the main function of the crest, and any other functions were secondary. Scientific interpretations of the crest's function began in 1910, when George Francis Eaton proposed two possibilities: an aerodynamic counterbalance and a muscle attachment point. He suggested that the crest might have anchored large, long jaw muscles, but admitted that this function alone could not explain the large size of some crests. Bennett (1992) agreed with Eaton's own assessment that the crest was too large and variable to have been a muscle attachment site. Eaton had suggested that a secondary function of the crest might have been as a counterbalance against the long beak, reducing the need for heavy neck muscles to control the orientation of the head. Wind tunnel tests showed that the crest did function as an effective counterbalance to a degree, but Bennett noted that, again, the hypothesis focuses only on the long crests of male P. longiceps, not on the larger crests of P. sternbergi and very small crests that existed among the females. Bennett found that the crests of females had no counterbalancing effect, and that the crests of male P. sternbergi would, by themselves, have a negative effect on the balance of the head. In fact, side to side movement of the crests would have required more, not less, neck musculature to control balance. In 1943, Dominik von Kripp suggested that the crest may have served as a rudder, an idea embraced by several later researchers. One researcher, Ross S. Stein, even suggested that the crest may have supported a membrane of skin connecting the backward-pointing crest to the neck and back, increasing its surface area and effectiveness as a rudder. The rudder hypothesis, again, does not take into account females nor P. sternbergi, which had an upward-pointing, not backward-pointing crest. Bennett also found that, even in its capacity as a rudder, the crest would not provide nearly so much directional force as simply maneuvering the wings. The suggestion that the crest was an air brake, and that the animals would turn their heads to the side in order to slow down, suffers from a similar problem. Additionally, the rudder and air brake hypotheses do not explain why such large variation exists in crest size even among adults. Alexander Kellner suggested that the large crests of the pterosaur Tapejara, as well as other species, might be used for heat exchange, allowing these pterosaurs to absorb or shed heat and regulate body temperature, which also would account for the correlation between crest size and body size. There is no evidence of extra blood vessels in the crest for this purpose, however, and the large, membranous wings filled with blood vessels would have served that purpose much more effectively. With these hypotheses ruled out, the best-supported hypothesis for crest function seems to be as a sexual display. This is consistent with the size variation seen in fossil specimens, where females and juveniles have small crests and males large, elaborate, variable crests. Sexual variation Adult Pteranodon specimens may be divided into two distinct size classes, small and large, with the large size class being about one and a half times larger than the small class, and the small class being twice as common as the large class. Both size classes lived alongside each other, and while researchers had previously suggested that they represent different species, Christopher Bennett showed that the differences between them are consistent with the concept that they represent females and males, and that Pteranodon species were sexually dimorphic. Skulls from the larger size class preserve large, upward and backward pointing crests, while the crests of the smaller size class are small and triangular. Some larger skulls also show evidence of a second crest that extended long and low, toward the tip of the beak, which is not seen in smaller specimens. The gender of the different size classes was determined, not from the skulls, but from the pelvic bones. Contrary to what may be expected, the smaller size class had disproportionately large and wide-set pelvic bones. Bennett interpreted this as indicating a more spacious birth canal, through which eggs would pass. He concluded that the small size class with small, triangular crests represent females, and the larger, large-crested specimens represent males. The overall size and crest size also corresponds to age. Immature specimens are known from both females and males, and immature males often have small crests similar to adult females. Therefore, it seems that the large crests only developed in males when they reached their large, adult size, making the sex of immature specimens difficult to establish from partial remains. The fact that females appear to have outnumbered males two to one suggests that, as with modern animals with size-related sexual dimorphism, such as sea lions and other pinnipeds, Pteranodon might have been polygynous, with a few males competing for association with groups consisting of large numbers of females. Similar to modern pinnipeds, Pteranodon may have competed to establish territory on rocky, offshore rookeries, with the largest, and largest-crested, males gaining the most territory and having more success mating with females. The crests of male Pteranodon would not have been used in competition, but rather as "visual dominance-rank symbols", with display rituals taking the place of physical competition with other males. If this hypothesis is correct, it also is likely that male Pteranodon played little to no part in rearing the young; such a behavior is not found in the males of modern polygynous animals who father many offspring at the same time. Paleoecology Specimens assigned to Pteranodon have been found in both the Smoky Hill Chalk deposits of the Niobrara Formation, and the slightly younger Sharon Springs deposits of the Pierre Shale Formation. When Pteranodon was alive, this area was covered by a large inland sea, known as the Western Interior Seaway. Famous for fossils collected since 1870, these formations extend from as far south as Kansas in the United States to Manitoba in Canada. However, Pteranodon specimens (or any pterosaur specimens) have only been found in the southern half of the formation, in Kansas, Wyoming, and South Dakota. Despite the fact that numerous fossils have been found in the contemporary parts of the formation in Canada, no pterosaur specimens have ever been found there. This strongly suggests that the natural geographic range of Pteranodon covered only the southern part of the Niobrara, and that its habitat did not extend farther north than South Dakota. Some very fragmentary fossils belonging to pteranodontian pterosaurs, and possibly Pteranodon itself, have also been found on the Gulf Coast and East Coast of the United States. For example, some bone fragments from the Mooreville Formation of Alabama and the Merchantville Formation of Delaware may have come from Pteranodon, though they are too incomplete to make a definite identification. Some remains from Japan have also been tentatively attributed to Pteranodon, but their distance from its known Western Interior Seaway habitat makes this identification unlikely. Pteranodon longiceps would have shared the sky with the giant-crested pterosaur Nyctosaurus. Compared to P. longiceps, which was a very common species, Nyctosaurus was rare, making up only 3% of pterosaur fossils from the formation. Also less common was the early toothed bird, Ichthyornis. Below the surface, the sea was populated primarily by invertebrates such as ammonites and squid. Vertebrate life, apart from basal fish, included sea turtles, such as Toxochelys, the plesiosaurs Elasmosaurus and Styxosaurus, and the flightless diving bird Parahesperornis. Mosasaurs were the most common marine reptiles, with genera including Clidastes, Mosasaurus and Tylosaurus. At least some of these marine reptiles are known to have fed on Pteranodon. Barnum Brown, in 1904, reported plesiosaur stomach contents containing "pterodactyl" bones, most likely from Pteranodon. Fossils from terrestrial dinosaurs also have been found in the Niobrara Chalk, suggesting that animals who died on shore must have been washed out to sea (one specimen of a hadrosaur appears to have been scavenged by a shark). It is likely that, as in other polygynous animals (in which males compete for association with harems of females), Pteranodon lived primarily on offshore rookeries, where they could nest away from land-based predators and feed far from shore; most Pteranodon fossils are found in locations which at the time, were hundreds of kilometres from the coastline. Classification Timespan and evolution Pteranodon fossils are known primarily from the Niobrara Formation of the central United States. Broadly defined, Pteranodon existed for more than four million years, during the Santonian stage of the Cretaceous period. The genus is present in most layers of the Niobrara Formation except for the upper two; in 2003, Kenneth Carpenter surveyed the distribution and dating of fossils in this formation, demonstrating that Pteranodon sternbergi existed there from 88 to 85 million years ago, while P. longiceps existed between 86 and 84.5 million years ago. A possible third species, which Kellner named Geosternbergia maiseyi in 2010, is known from the Sharon Springs member of the Pierre Shale Formation in Kansas, Wyoming, and South Dakota, dating to between 81.5 and 80.5 million years ago. In the early 1990s, Bennett noted that the two major morphs of pteranodont present in the Niobrara Formation were precisely separated in time with little, if any, overlap. Due to this, and to their gross overall similarity, he suggested that they probably represent chronospecies within a single evolutionary lineage lasting about 4 million years. In other words, only one species of Pteranodon would have been present at any one time, and P. sternbergi (or Geosternbergia) in all likelihood was the direct ancestor species of P. longiceps. Valid species Many researchers consider there to be at least two species of Pteranodon. However, aside from the differences between males and females described above, the post-cranial skeletons of Pteranodon show little to no variation between species or specimens, and the bodies and wings of all pteranodonts were essentially identical. Two species of Pteranodon are traditionally recognized as valid: Pteranodon longiceps, the type species, and Pteranodon sternbergi. The species differ only in the shape of the crest in adult males (described above), and possibly in the angle of certain skull bones. Because well-preserved Pteranodon skull fossils are extremely rare, researchers use stratigraphy (i.e. which rock layer of the geologic formation a fossil is found in) to determine species identity in most cases. Pteranodon sternbergi is the only known species of Pteranodon with an upright crest. The lower jaw of P. sternbergi was long. It was collected by George F. Sternberg in 1952 and described by John Christian Harksen in 1966, from the lower portion of the Niobrara Formation. It was older than P. longiceps and is considered by Bennett to be the direct ancestor of the later species. Because fossils identifiable as P. sternbergi are found exclusively in the lower layers of the Niobrara Formation, and P. longiceps fossils exclusively in the upper layers, a fossil lacking the skull can be identified based on its position in the geologic column (though for many early fossil finds, precise data about its location was not recorded, rendering many fossils unidentifiable). Below is a cladogram showing the phylogenetic placement of this genus within Pteranodontia from Andres and Myers (2013). Alternative classifications Due to the subtle variations between specimens of pteranodontid from the Niobrara Formation, most researchers have assigned all of them to the single genus Pteranodon, in at least two species (P. longiceps and P. sternbergi) distinguished mainly by the shape of the crest. However, the classification of these two forms has varied from researcher to researcher. In 1972, Halsey Wilkinson Miller published a paper arguing that the various forms of Pteranodon were different enough to be placed in distinct subgenera. He named these Pteranodon (Occidentalia) occidentalis (for the now-disused species P. occidentalis) and Pteranodon (Sternbergia) sternbergi. However, the name Sternbergia was preoccupied, and in 1978 Miller re-named the species Pteranodon (Geosternbergia) sternbergi, and named a third subgenus/species combination for P. longiceps, as Pteranodon (Longicepia) longiceps. Most prominent pterosaur researchers of the late 20th century however, including S. Christopher Bennett and Peter Wellnhofer, did not adopt these subgeneric names, and continued to place all pteranodont species into the single genus Pteranodon. In 2010, pterosaur researcher Alexander Kellner revisited H.W. Miller's classification. Kellner followed Miller's opinion that the differences between the Pteranodon species were great enough to place them into different genera. He placed P. sternbergi into the genus named by Miller, Geosternbergia, along with the Pierre Shale skull specimen which Bennett had previously considered to be a large male P. longiceps. Kellner argued that this specimen's crest, though incompletely preserved, was most similar to Geosternbergia. Because the specimen was millions of years younger than any known Geosternbergia, he assigned it to the new species Geosternbergia maiseyi. Numerous other pteranodont specimens are known from the same formation and time period, and Kellner suggested they may belong to the same species as G. maiseyi, but because they lack skulls, he could not confidently identify them. However, both species previously referred to Geosternbergia were separately included as those of Pteranodon (P. sternbergi and P. maiseyi) based on phylogenetic analysis in 2024. Disused species A number of additional species of Pteranodon have been named since the 1870s, although most now are considered to be junior synonyms of two or three valid species. The best-supported is the type species, P. longiceps, based on the well-preserved specimen including the first-known skull found by S. W. Williston. This individual had a wingspan of . Other valid species include the possibly larger P. sternbergi, with a wingspan originally estimated at . P. oweni (P. occidentalis), P. velox, P. umbrosus, P. harpyia, and P. comptus are considered to be nomina dubia by Bennett (1994) and others who question their validity. All probably are synonymous with the more well-known species. Because the key distinguishing characteristic Marsh noted for Pteranodon was its lack of teeth, any toothless pterosaur jaw fragment, wherever it was found in the world, tended to be attributed to Pteranodon during the late nineteenth and early twentieth centuries. This resulted in a plethora of species and a great deal of confusion. The name became a wastebasket taxon, rather like the dinosaur Megalosaurus, to label any pterosaur remains that could not be distinguished other than by the absence of teeth. Species (often dubious ones now known to be based on sexual variation or juvenile characters) have been reclassified a number of times, and several subgenera have in the 1970s been erected by Halsey Wilkinson Miller to hold them in various combinations, further confusing the taxonomy (subgenera include Longicepia, Occidentalia, and Geosternbergia). Notable authors who have discussed the various aspects of Pteranodon include Bennett, Padian, Unwin, Kellner, and Wellnhofer. Two species, P. oregonensis and P. orientalis, are not pteranodontids and have been renamed Bennettazhia oregonensis and Bogolubovia orientalis respectively. List of species and synonyms Status of names listed below follow a survey by Bennett, 1994 unless otherwise noted.	Biology and health sciences	Pterosaurs	Animals
24829	https://en.wikipedia.org/wiki/Primitive%20recursive%20function	Primitive recursive function	In computability theory, a primitive recursive function is, roughly speaking, a function that can be computed by a computer program whose loops are all "for" loops (that is, an upper bound of the number of iterations of every loop is fixed before entering the loop). Primitive recursive functions form a strict subset of those general recursive functions that are also total functions. The importance of primitive recursive functions lies in the fact that most computable functions that are studied in number theory (and more generally in mathematics) are primitive recursive. For example, addition and division, the factorial and exponential function, and the function which returns the nth prime are all primitive recursive. In fact, for showing that a computable function is primitive recursive, it suffices to show that its time complexity is bounded above by a primitive recursive function of the input size. It is hence not particularly easy to devise a computable function that is not primitive recursive; some examples are shown in section below. The set of primitive recursive functions is known as PR in computational complexity theory. Definition A primitive recursive function takes a fixed number of arguments, each a natural number (nonnegative integer: {0, 1, 2, ...}), and returns a natural number. If it takes n arguments it is called n-ary. The basic primitive recursive functions are given by these axioms: More complex primitive recursive functions can be obtained by applying the operations given by these axioms: The primitive recursive functions are the basic functions and those obtained from the basic functions by applying these operations a finite number of times. Examples Addition A definition of the 2-ary function , to compute the sum of its arguments, can be obtained using the primitive recursion operator . To this end, the well-known equations are "rephrased in primitive recursive function terminology": In the definition of , the first equation suggests to choose to obtain ; the second equation suggests to choose to obtain . Therefore, the addition function can be defined as . As a computation example, Doubling Given , the 1-ary function doubles its argument, . Multiplication In a similar way as addition, multiplication can be defined by . This reproduces the well-known multiplication equations: and Predecessor The predecessor function acts as the "opposite" of the successor function and is recursively defined by the rules and . A primitive recursive definition is . As a computation example, Truncated subtraction The limited subtraction function (also called "monus", and denoted "") is definable from the predecessor function. It satisfies the equations Since the recursion runs over the second argument, we begin with a primitive recursive definition of the reversed subtraction, . Its recursion then runs over the first argument, so its primitive recursive definition can be obtained, similar to addition, as . To get rid of the reversed argument order, then define . As a computation example, Converting predicates to numeric functions In some settings it is natural to consider primitive recursive functions that take as inputs tuples that mix numbers with truth values (that is for true and for false), or that produce truth values as outputs. This can be accomplished by identifying the truth values with numbers in any fixed manner. For example, it is common to identify the truth value with the number and the truth value with the number . Once this identification has been made, the characteristic function of a set , which always returns or , can be viewed as a predicate that tells whether a number is in the set . Such an identification of predicates with numeric functions will be assumed for the remainder of this article. Predicate "Is zero" As an example for a primitive recursive predicate, the 1-ary function shall be defined such that if , and , otherwise. This can be achieved by defining . Then, and e.g. . Predicate "Less or equal" Using the property , the 2-ary function can be defined by . Then if , and , otherwise. As a computation example, Predicate "Greater or equal" Once a definition of is obtained, the converse predicate can be defined as . Then, is true (more precisely: has value 1) if, and only if, . If-then-else The 3-ary if-then-else operator known from programming languages can be defined by . Then, for arbitrary , and . That is, returns the then-part, , if the if-part, , is true, and the else-part, , otherwise. Junctors Based on the function, it is easy to define logical junctors. For example, defining , one obtains , that is, is true if, and only if, both and are true (logical conjunction of and ). Similarly, and lead to appropriate definitions of disjunction and negation: and . Equality predicate Using the above functions , and , the definition implements the equality predicate. In fact, is true if, and only if, equals . Similarly, the definition implements the predicate "less-than", and implements "greater-than". Other operations on natural numbers Exponentiation and primality testing are primitive recursive. Given primitive recursive functions , , , and , a function that returns the value of when and the value of otherwise is primitive recursive. Operations on integers and rational numbers By using Gödel numberings, the primitive recursive functions can be extended to operate on other objects such as integers and rational numbers. If integers are encoded by Gödel numbers in a standard way, the arithmetic operations including addition, subtraction, and multiplication are all primitive recursive. Similarly, if the rationals are represented by Gödel numbers then the field operations are all primitive recursive. Some common primitive recursive functions The following examples and definitions are from . Many appear with proofs. Most also appear with similar names, either as proofs or as examples, in they add the logarithm lo(x, y) or lg(x, y) depending on the exact derivation. In the following the mark " ' ", e.g. a', is the primitive mark meaning "the successor of", usually thought of as " +1", e.g. a +1 =def a'. The functions 16–20 and #G are of particular interest with respect to converting primitive recursive predicates to, and extracting them from, their "arithmetical" form expressed as Gödel numbers. Addition: a+b Multiplication: a×b Exponentiation: ab Factorial a! : 0! = 1, a'! = a!×a' pred(a): (Predecessor or decrement): If a > 0 then a−1 else 0 Proper subtraction a ∸ b: If a ≥ b then a−b else 0 Minimum(a1, ... an) Maximum(a1, ... an) Absolute difference: \| a−b \| =def (a ∸ b) + (b ∸ a) ~sg(a): NOT[signum(a)]: If a=0 then 1 else 0 sg(a): signum(a): If a=0 then 0 else 1 a \| b: (a divides b): If b=k×a for some k then 0 else 1 Remainder(a, b): the leftover if b does not divide a "evenly". Also called MOD(a, b) a = b: sg \| a − b \| (Kleene's convention was to represent true by 0 and false by 1; presently, especially in computers, the most common convention is the reverse, namely to represent true by 1 and false by 0, which amounts to changing sg into ~sg here and in the next item) a < b: sg( a' ∸ b ) Pr(a): a is a prime number Pr(a) =def a>1 & NOT(Exists c)1<c<a [ c\|a ] pi: the i+1th prime number (a)i: exponent of pi in a: the unique x such that pix\|a & NOT(pix'\|a) lh(a): the "length" or number of non-vanishing exponents in a lo(a, b): (logarithm of a to base b): If a, b > 1 then the greatest x such that bx \| a else 0 In the following, the abbreviation x =def x1, ... xn; subscripts may be applied if the meaning requires. #A: A function φ definable explicitly from functions Ψ and constants q1, ... qn is primitive recursive in Ψ. #B: The finite sum Σy<z ψ(x, y) and product Πy<zψ(x, y) are primitive recursive in ψ. #C: A predicate P obtained by substituting functions χ1,..., χm for the respective variables of a predicate Q is primitive recursive in χ1,..., χm, Q. #D: The following predicates are primitive recursive in Q and R: NOT_Q(x) . Q OR R: Q(x) V R(x), Q AND R: Q(x) & R(x), Q IMPLIES R: Q(x) → R(x) Q is equivalent to R: Q(x) ≡ R(x) #E: The following predicates are primitive recursive in the predicate R: (Ey)y<z R(x, y) where (Ey)y<z denotes "there exists at least one y that is less than z such that" (y)y<z R(x, y) where (y)y<z denotes "for all y less than z it is true that" μyy<z R(x, y). The operator μyy<z R(x, y) is a bounded form of the so-called minimization- or mu-operator: Defined as "the least value of y less than z such that R(x, y) is true; or z if there is no such value." #F: Definition by cases: The function defined thus, where Q1, ..., Qm are mutually exclusive predicates (or "ψ(x) shall have the value given by the first clause that applies), is primitive recursive in φ1, ..., Q1, ... Qm: φ(x) = φ1(x) if Q1(x) is true, . . . . . . . . . . . . . . . . . . . φm(x) if Qm(x) is true φm+1(x) otherwise #G: If φ satisfies the equation: φ(y,x) = χ(y, COURSE-φ(y; x2, ... xn ), x2, ... xn then φ is primitive recursive in χ. The value COURSE-φ(y; x2 to n ) of the course-of-values function encodes the sequence of values φ(0,x2 to n), ..., φ(y-1,x2 to n) of the original function. Relationship to recursive functions The broader class of partial recursive functions is defined by introducing an unbounded search operator. The use of this operator may result in a partial function, that is, a relation with at most one value for each argument, but does not necessarily have any value for any argument (see domain). An equivalent definition states that a partial recursive function is one that can be computed by a Turing machine. A total recursive function is a partial recursive function that is defined for every input. Every primitive recursive function is total recursive, but not all total recursive functions are primitive recursive. The Ackermann function A(m,n) is a well-known example of a total recursive function (in fact, provable total), that is not primitive recursive. There is a characterization of the primitive recursive functions as a subset of the total recursive functions using the Ackermann function. This characterization states that a function is primitive recursive if and only if there is a natural number m such that the function can be computed by a Turing machine that always halts within A(m,n) or fewer steps, where n is the sum of the arguments of the primitive recursive function. An important property of the primitive recursive functions is that they are a recursively enumerable subset of the set of all total recursive functions (which is not itself recursively enumerable). This means that there is a single computable function f(m,n) that enumerates the primitive recursive functions, namely: For every primitive recursive function g, there is an m such that g(n) = f(m,n) for all n, and For every m, the function h(n) = f(m,n) is primitive recursive. f can be explicitly constructed by iteratively repeating all possible ways of creating primitive recursive functions. Thus, it is provably total. One can use a diagonalization argument to show that f is not recursive primitive in itself: had it been such, so would be h(n) = f(n,n)+1. But if this equals some primitive recursive function, there is an m such that h(n) = f(m,n) for all n, and then h(m) = f(m,m), leading to contradiction. However, the set of primitive recursive functions is not the largest recursively enumerable subset of the set of all total recursive functions. For example, the set of provably total functions (in Peano arithmetic) is also recursively enumerable, as one can enumerate all the proofs of the theory. While all primitive recursive functions are provably total, the converse is not true. Limitations Primitive recursive functions tend to correspond very closely with our intuition of what a computable function must be. Certainly the initial functions are intuitively computable (in their very simplicity), and the two operations by which one can create new primitive recursive functions are also very straightforward. However, the set of primitive recursive functions does not include every possible total computable function—this can be seen with a variant of Cantor's diagonal argument. This argument provides a total computable function that is not primitive recursive. A sketch of the proof is as follows: This argument can be applied to any class of computable (total) functions that can be enumerated in this way, as explained in the article Machine that always halts. Note however that the partial computable functions (those that need not be defined for all arguments) can be explicitly enumerated, for instance by enumerating Turing machine encodings. Other examples of total recursive but not primitive recursive functions are known: The function that takes m to Ackermann(m,m) is a unary total recursive function that is not primitive recursive. The Paris–Harrington theorem involves a total recursive function that is not primitive recursive. The Sudan function The Goodstein function Variants Constant functions Instead of , alternative definitions use just one 0-ary zero function as a primitive function that always returns zero, and built the constant functions from the zero function, the successor function and the composition operator. Iterative functions Robinson considered various restrictions of the recursion rule. One is the so-called iteration rule where the function h does not have access to the parameters xi (in this case, we may assume without loss of generality that the function g is just the identity, as the general case can be obtained by substitution): He proved that the class of all primitive recursive functions can still be obtained in this way. Pure recursion Another restriction considered by Robinson is pure recursion, where h does not have access to the induction variable y: Gladstone proved that this rule is enough to generate all primitive recursive functions. Gladstone improved this so that even the combination of these two restrictions, i.e., the pure iteration rule below, is enough: Further improvements are possible: Severin prove that even the pure iteration rule without parameters, namely suffices to generate all unary primitive recursive functions if we extend the set of initial functions with truncated subtraction x ∸ y. We get all primitive recursive functions if we additionally include + as an initial function. Additional primitive recursive forms Some additional forms of recursion also define functions that are in fact primitive recursive. Definitions in these forms may be easier to find or more natural for reading or writing. Course-of-values recursion defines primitive recursive functions. Some forms of mutual recursion also define primitive recursive functions. The functions that can be programmed in the LOOP programming language are exactly the primitive recursive functions. This gives a different characterization of the power of these functions. The main limitation of the LOOP language, compared to a Turing-complete language, is that in the LOOP language the number of times that each loop will run is specified before the loop begins to run. Computer language definition An example of a primitive recursive programming language is one that contains basic arithmetic operators (e.g. + and −, or ADD and SUBTRACT), conditionals and comparison (IF-THEN, EQUALS, LESS-THAN), and bounded loops, such as the basic for loop, where there is a known or calculable upper bound to all loops (FOR i FROM 1 TO n, with neither i nor n modifiable by the loop body). No control structures of greater generality, such as while loops or IF-THEN plus GOTO, are admitted in a primitive recursive language. The LOOP language, introduced in a 1967 paper by Albert R. Meyer and Dennis M. Ritchie, is such a language. Its computing power coincides with the primitive recursive functions. A variant of the LOOP language is Douglas Hofstadter's BlooP in Gödel, Escher, Bach. Adding unbounded loops (WHILE, GOTO) makes the language general recursive and Turing-complete, as are all real-world computer programming languages. The definition of primitive recursive functions implies that their computation halts on every input (after a finite number of steps). On the other hand, the halting problem is undecidable for general recursive functions. Finitism and consistency results The primitive recursive functions are closely related to mathematical finitism, and are used in several contexts in mathematical logic where a particularly constructive system is desired. Primitive recursive arithmetic (PRA), a formal axiom system for the natural numbers and the primitive recursive functions on them, is often used for this purpose. PRA is much weaker than Peano arithmetic, which is not a finitistic system. Nevertheless, many results in number theory and in proof theory can be proved in PRA. For example, Gödel's incompleteness theorem can be formalized into PRA, giving the following theorem: If T is a theory of arithmetic satisfying certain hypotheses, with Gödel sentence GT, then PRA proves the implication Con(T)→GT. Similarly, many of the syntactic results in proof theory can be proved in PRA, which implies that there are primitive recursive functions that carry out the corresponding syntactic transformations of proofs. In proof theory and set theory, there is an interest in finitistic consistency proofs, that is, consistency proofs that themselves are finitistically acceptable. Such a proof establishes that the consistency of a theory T implies the consistency of a theory S by producing a primitive recursive function that can transform any proof of an inconsistency from S into a proof of an inconsistency from T. One sufficient condition for a consistency proof to be finitistic is the ability to formalize it in PRA. For example, many consistency results in set theory that are obtained by forcing can be recast as syntactic proofs that can be formalized in PRA. History Recursive definitions had been used more or less formally in mathematics before, but the construction of primitive recursion is traced back to Richard Dedekind's theorem 126 of his Was sind und was sollen die Zahlen? (1888). This work was the first to give a proof that a certain recursive construction defines a unique function. Primitive recursive arithmetic was first proposed by Thoralf Skolem in 1923. The current terminology was coined by Rózsa Péter (1934) after Ackermann had proved in 1928 that the function which today is named after him was not primitive recursive, an event which prompted the need to rename what until then were simply called recursive functions.	Mathematics	Computability theory	null
24834	https://en.wikipedia.org/wiki/Protein%20targeting	Protein targeting	Protein targeting or protein sorting is the biological mechanism by which proteins are transported to their appropriate destinations within or outside the cell. Proteins can be targeted to the inner space of an organelle, different intracellular membranes, the plasma membrane, or to the exterior of the cell via secretion. Information contained in the protein itself directs this delivery process. Correct sorting is crucial for the cell; errors or dysfunction in sorting have been linked to multiple diseases. History In 1970, Günter Blobel conducted experiments on protein translocation across membranes. Blobel, then an assistant professor at Rockefeller University, built upon the work of his colleague George Palade. Palade had previously demonstrated that non-secreted proteins were translated by free ribosomes in the cytosol, while secreted proteins (and target proteins, in general) were translated by ribosomes bound to the endoplasmic reticulum. Candidate explanations at the time postulated a processing difference between free and ER-bound ribosomes, but Blobel hypothesized that protein targeting relied on characteristics inherent to the proteins, rather than a difference in ribosomes. Supporting his hypothesis, Blobel discovered that many proteins have a short amino acid sequence at one end that functions like a postal code specifying an intracellular or extracellular destination. He described these short sequences (generally 13 to 36 amino acids residues) as signal peptides or signal sequences and was awarded the 1999 Nobel prize in Physiology for the same. Signal peptides Signal peptides serve as targeting signals, enabling cellular transport machinery to direct proteins to specific intracellular or extracellular locations. While no consensus sequence has been identified for signal peptides, many nonetheless possess a characteristic tripartite structure: A positively charged, hydrophilic region near the N-terminal. A span of 10 to 15 hydrophobic amino acids near the middle of the signal peptide. A slightly polar region near the C-terminal, typically favoring amino acids with smaller side chains at positions approaching the cleavage site. After a protein has reached its destination, the signal peptide is generally cleaved by a signal peptidase. Consequently, most mature proteins do not contain signal peptides. While most signal peptides are found at the N-terminal, in peroxisomes the targeting sequence is located on the C-terminal extension. Unlike signal peptides, signal patches are composed by amino acid residues that are discontinuous in the primary sequence but become functional when folding brings them together on the protein surface. Unlike most signal sequences, signal patches are not cleaved after sorting is complete. In addition to intrinsic signaling sequences, protein modifications like glycosylation can also induce targeting to specific intracellular or extracellular regions. Protein translocation Since the translation of mRNA into protein by a ribosome takes place within the cytosol, proteins destined for secretion or a specific organelle must be translocated. This process can occur during translation, known as co-translational translocation, or after translation is complete, known as post-translational translocation. Co-translational translocation Most secretory and membrane-bound proteins are co-translationally translocated. Proteins that reside in the endoplasmic reticulum (ER), golgi or endosomes also use the co-translational translocation pathway. This process begins while the protein is being synthesized on the ribosome, when a signal recognition particle (SRP) recognizes an N-terminal signal peptide of the nascent protein. Binding of the SRP temporarily pauses synthesis while the ribosome-protein complex is transferred to an SRP receptor on the ER in eukaryotes, and the plasma membrane in prokaryotes. There, the nascent protein is inserted into the translocon, a membrane-bound protein conducting channel composed of the Sec61 translocation complex in eukaryotes, and the homologous SecYEG complex in prokaryotes. In secretory proteins and type I transmembrane proteins, the signal sequence is immediately cleaved from the nascent polypeptide once it has been translocated into the membrane of the ER (eukaryotes) or plasma membrane (prokaryotes) by signal peptidase. The signal sequence of type II membrane proteins and some polytopic membrane proteins are not cleaved off and therefore are referred to as signal anchor sequences. Within the ER, the protein is first covered by a chaperone protein to protect it from the high concentration of other proteins in the ER, giving it time to fold correctly. Once folded, the protein is modified as needed (for example, by glycosylation), then transported to the Golgi for further processing and goes to its target organelles or is retained in the ER by various ER retention mechanisms. The amino acid chain of transmembrane proteins, which often are transmembrane receptors, passes through a membrane one or several times. These proteins are inserted into the membrane by translocation, until the process is interrupted by a stop-transfer sequence, also called a membrane anchor or signal-anchor sequence. These complex membrane proteins are currently characterized using the same model of targeting that has been developed for secretory proteins. However, many complex multi-transmembrane proteins contain structural aspects that do not fit this model. Seven transmembrane G-protein coupled receptors (which represent about 5% of the genes in humans) mostly do not have an amino-terminal signal sequence. In contrast to secretory proteins, the first transmembrane domain acts as the first signal sequence, which targets them to the ER membrane. This also results in the translocation of the amino terminus of the protein into the ER membrane lumen. This translocation, which has been demonstrated with opsin with in vitro experiments, breaks the usual pattern of "co-translational" translocation which has always held for mammalian proteins targeted to the ER. A great deal of the mechanics of transmembrane topology and folding remains to be elucidated. Post-translational translocation Even though most secretory proteins are co-translationally translocated, some are translated in the cytosol and later transported to the ER/plasma membrane by a post-translational system. In prokaryotes this process requires certain cofactors such as SecA and SecB and is facilitated by Sec62 and Sec63, two membrane-bound proteins. The Sec63 complex, which is embedded in the ER membrane, causes hydrolysis of ATP, allowing chaperone proteins to bind to an exposed peptide chain and slide the polypeptide into the ER lumen. Once in the lumen the polypeptide chain can be folded properly. This process only occurs in unfolded proteins located in the cytosol. In addition, proteins targeted to other cellular destinations, such as mitochondria, chloroplasts, or peroxisomes, use specialized post-translational pathways. Proteins targeted for the nucleus are also translocated post-translationally through the addition of a nuclear localization signal (NLS) that promotes passage through the nuclear envelope via nuclear pores. Sorting of proteins Mitochondria While some proteins in the mitochondria originate from mitochondrial DNA within the organelle, most mitochondrial proteins are synthesized as cytosolic precursors containing uptake peptide signals. Unfolded proteins bound by cytosolic chaperone hsp70 that are targeted to the mitochondria may be localized to four different areas depending on their sequences. They may be targeted to the mitochondrial matrix, the outer membrane, the intermembrane space, or the inner membrane. Defects in any one or more of these processes has been linked to health and disease. Mitochondrial matrix Proteins destined for the mitochondrial matrix have specific signal sequences at their beginning (N-terminus) that consist of a string of 20 to 50 amino acids. These sequences are designed to interact with receptors that guide the proteins to their correct location inside the mitochondria. The sequences have a unique structure with clusters of water-loving (hydrophilic) and water-avoiding (hydrophobic) amino acids, giving them a dual nature known as amphipathic. These amphipathic sequences typically form a spiral shape (alpha-helix) with the charged amino acids on one side and the hydrophobic ones on the opposite side. This structural feature is essential for the sequence to function correctly in directing proteins to the matrix. If mutations occur that mess with this dual nature, the protein often fails to reach its intended destination, although not all changes to the sequence have this effect. This indicates the importance of the amphipathic property for the protein to be correctly targeted to the mitochondrial matrix.Proteins targeted to the mitochondrial matrix first involves interactions between the matrix targeting sequence located at the N-terminus and the outer membrane import receptor complex TOM20/22. In addition to the docking of internal sequences and cytosolic chaperones to TOM70. Where TOM is an abbreviation for translocase of the outer membrane. Binding of the matrix targeting sequence to the import receptor triggers a handoff of the polypeptide to the general import core (GIP) known as TOM40. The general import core (TOM40) then feeds the polypeptide chain through the intermembrane space and into another translocase complex TIM17/23/44 located on the inner mitochondrial membrane. This is accompanied by the necessary release of the cytosolic chaperones that maintain an unfolded state prior to entering the mitochondria. As the polypeptide enters the matrix, the signal sequence is cleaved by a processing peptidase and the remaining sequences are bound by mitochondrial chaperones to await proper folding and activity. The push and pull of the polypeptide from the cytosol to the intermembrane space and then the matrix is achieved by an electrochemical gradient that is established by the mitochondrion during oxidative phosphorylation. In which a mitochondrion active in metabolism has generated a negative potential inside the matrix and a positive potential in the intermembrane space. It is this negative potential inside the matrix that directs the positively charged regions of the targeting sequence into its desired location. Mitochondrial inner membrane Targeting of mitochondrial proteins to the inner membrane may follow 3 different pathways depending upon their overall sequences, however, entry from the outer membrane remains the same using the import receptor complex TOM20/22 and TOM40 general import core. The first pathway for proteins targeted to the inner membrane follows the same steps as those designated to the matrix where it contains a matrix targeting sequence that channels the polypeptide to the inner membrane complex containing the previously mentioned translocase complex TIM17/23/44. However, the difference is that the peptides that are designated to the inner membrane and not the matrix contain an upstream sequence called the stop-transfer-anchor sequence. This stop-transfer-anchor sequence is a hydrophobic region that embeds itself into the phospholipid bilayer of the inner membrane and prevents translocation further into the mitochondrion. The second pathway for proteins targeted to the inner membrane follows the matrix localization pathway in its entirety. However, instead of a stop-transfer-anchor sequence, it contains another sequence that interacts with an inner membrane protein called Oxa-1 once inside the matrix that will embed it into the inner membrane. The third pathway for mitochondrial proteins targeted to the inner membrane follow the same entry as the others into the outer membrane, however, this pathway utilizes the translocase complex TIM22/54 assisted by complex TIM9/10 in the intermembrane space to anchor the incoming peptide into the membrane. The peptides for this last pathway do not contain a matrix targeting sequence, but instead contain several internal targeting sequences. Mitochondrial intermembrane space If instead the precursor protein is designated to the intermembrane space of the mitochondrion, there are two pathways this may occur depending on the sequences being recognized. The first pathway to the intermembrane space follows the same steps for an inner membrane targeted protein. However, once bound to the inner membrane the C-terminus of the anchored protein is cleaved via a peptidase that liberates the preprotein into the intermembrane space so it can fold into its active state. One of the greatest examples for a protein that follows this pathway is cytochrome b2, that upon being cleaved will interact with a heme cofactor and become active. The second intermembrane space pathway does not utilize any inner membrane complexes and therefor does not contain a matrix targeting signal. Instead, it enters through the general import core TOM40 and is further modified in the intermembrane space to achieve its active conformation. TIM9/10 is an example of a protein that follows this pathway in order to be in the location it needs to be to assist in inner membrane targeting. Mitochondrial outer membrane Outer membrane targeting simply involves the interaction of precursor proteins with the outer membrane translocase complexes that embeds it into the membrane via internal-targeting sequences that are to form hydrophobic alpha helices or beta barrels that span the phospholipid bilayer. This may occur by two different routes depending on the preprotein internal sequences. If the preprotein contains internal hydrophobic regions capable of forming alpha helices, then the preprotein will utilize the mitochondrial import complex (MIM) and be transferred laterally to the membrane. For preproteins containing hydrophobic internal sequences that correlate to beta-barrel forming proteins, they will be imported from the aforementioned outer membrane complex TOM20/22 to the intermembrane space. In which they will interact with TIM9/10 intermembrane-space protein complex that transfers them to sorting and assembly machinery (SAM) that is present in the outer membrane that laterally displaces the targeted protein as a beta-barrel. Chloroplasts Chloroplasts are similar to mitochondria in that they contain their own DNA for production of some of their components. However, the majority of their proteins are obtained via post-translational translocation and arise from nuclear genes. Proteins may be targeted to several sites of the chloroplast depending on their sequences such as the outer envelope, inner envelope, stroma, thylakoid lumen, or the thylakoid membrane. Proteins are targeted to Thylakoids by mechanisms related to Bacterial Protein Translocation. Proteins targeted to the envelope of chloroplasts usually lack cleavable sorting sequence and are laterally displaced via membrane sorting complexes. General import for the majority of preproteins requires translocation from the cytosol through the Toc and Tic complexes located within the chloroplast envelope. Where Toc is an abbreviation for the translocase of the outer chloroplast envelope and Tic is the translocase of the inner chloroplast envelope. There is a minimum of three proteins that make up the function of the Toc complex. Two of which, referred to as Toc159 and Toc34, are responsible for the docking of stromal import sequences and both contain GTPase activity. The third known as Toc 75, is the actual translocation channel that feeds the recognized preprotein by Toc159/34 into the chloroplast. Stroma Targeting to the stroma requires the preprotein to have a stromal import sequence that is recognized by the Tic complex of the inner envelope upon being translocated from the outer envelope by the Toc complex. The Tic complex is composed of at least five different Tic proteins that are required to form the translocation channel across the inner envelope. Upon being delivered to the stroma, the stromal import sequence is cleaved off via a signal peptidase. This delivery process to the stroma is currently known to be driven by ATP hydrolysis via stromal HSP chaperones, instead of the transmembrane electrochemical gradient that is established in mitochondria to drive protein import. Further intra-chloroplast sorting depends on additional target sequences such as those designated to the thylakoid membrane or the thylakoid lumen. Thylakoid lumen If a protein is to be targeted to the thylakoid lumen, this may occur via four differently known routes that closely resemble bacterial protein transport mechanisms. The route that is taken depends upon the protein delivered to the stroma being in either an unfolded or metal-bound folded state. Both of which will still contain a thylakoid targeting sequence that is also cleaved upon entry to the lumen. While protein import into the stroma is ATP-driven, the pathway for metal-bound proteins in a folded state to the thylakoid lumen has been shown to be driven by a pH gradient. Thylakoid membrane Proteins bound for the membrane of the thylakoid will follow up to four known routes that are illustrated in the corresponding figure shown. They may follow a co-translational insertion route that utilizes stromal ribosomes and the SecY/E transmembrane complex, the SRP-dependent pathway, the spontaneous insertion pathway, or the GET pathway. The last of the three are post-translational pathways originating from nuclear genes and therefor constitute the majority of proteins targeted to the thylakoid membrane. According to recent review articles in the journal of biochemistry and molecular biology, the exact mechanisms are not yet fully understood. Both chloroplasts and mitochondria Many proteins are needed in both mitochondria and chloroplasts. In general the dual-targeting peptide is of intermediate character to the two specific ones. The targeting peptides of these proteins have a high content of basic and hydrophobic amino acids, a low content of negatively charged amino acids. They have a lower content of alanine and a higher content of leucine and phenylalanine. The dual targeted proteins have a more hydrophobic targeting peptide than both mitochondrial and chloroplastic ones. However, it is tedious to predict if a peptide is dual-targeted or not based on its physio-chemical characteristics. Nucleus The nucleus of a cell is surrounded by a nuclear envelope consisting of two layers, with the inner layer providing structural support and anchorage for chromosomes and the nuclear lamina. The outer layer is similar to the endoplasmic reticulum (ER) membrane. This envelope contains nuclear pores, which are complex structures made from around 30 different proteins. These pores act as selective gates that control the flow of molecules into and out of the nucleus. While small molecules can pass through these pores without issue, larger molecules, like RNA and proteins destined for the nucleus, must have specific signals to be allowed through. These signals are known as nuclear localization signals, usually comprising short sequences rich in positively charged amino acids like lysine or arginine. Proteins called nuclear import receptors recognize these signals and guide the large molecules through the nuclear pores by interacting with the disordered, mesh-like proteins that fill the pore. The process is dynamic, with the receptor moving the molecule through the meshwork until it reaches the nucleus. Once inside, a GTPase enzyme called Ran, which can exist in two different forms (one bound to GTP and the other to GDP), facilitates the release of the cargo inside the nucleus and recycles the receptor back to the cytosol. The energy for this transport comes from the hydrolysis of GTP by Ran. Similarly, nuclear export receptors help move proteins and RNA out of the nucleus using a different signal and also harnessing Ran's energy conversion. Overall, the nuclear pore complex works efficiently to transport macromolecules at high speed, allowing proteins to move in their folded state and ribosomal components as complete particles, which is distinct from how proteins are transported into most other organelles. Endoplasmic reticulum The endoplasmic reticulum (ER) plays a key role in protein synthesis and distribution in eukaryotic cells. It's a vast network of membranes where proteins are processed and sorted to various destinations, including the ER itself, the cell surface, and other organelles like the Golgi apparatus, endosomes, and lysosomes. Unlike other organelle-targeted proteins, those headed for the ER start to be transferred across its membrane while they're still being made. Protein synthesis and sorting There are two types of proteins that move to the ER: water-soluble proteins, which completely cross into the ER lumen, and transmembrane proteins, which partly cross and embed themselves within the ER membrane. These proteins find their way to the ER with the help of an ER signal sequence, a short stretch of hydrophobic amino acids. Proteins entering the ER are synthesized by ribosomes. There are two sets of ribosomes in the cell: those bound to the ER (making it look 'rough') and those floating freely in the cytosol. Both sets are identical but differ in the proteins they synthesize at a given moment. Ribosomes that are making proteins with an ER signal sequence attach to the ER membrane and start the translocation process. This process is energy-efficient because the growing protein chain itself pushes through the ER membrane as it elongates. As the mRNA is translated into a protein, multiple ribosomes may attach to it, creating a structure called a polyribosome. If the mRNA is coding for a protein with an ER signal sequence, the polyribosome attaches to the ER membrane, and the protein begins to enter the ER while it is still being synthesized. Guided entry of soluble proteins In the process of protein synthesis within eukaryotic cells, soluble proteins that are destined for the endoplasmic reticulum (ER) or for secretion out of the cell are guided to the ER by a two-part system. Firstly, a signal-recognition particle (SRP) in the cytosol attaches to the emerging protein's ER signal sequence and the ribosome itself. Secondly, an SRP receptor located in the ER membrane recognizes and binds to the SRP. This interaction temporarily slows down protein synthesis until the SRP and ribs complex binds to the SRP receptor on the ER. Once this binding occurs, the SRP is released, and the ribosome is transferred to a protein translocator in the ER membrane, allowing protein synthesis to continue. The polypeptide chain of the protein is then threaded through a channel in the translocator into the ER lumen. The signal sequence of the protein, typically at the beginning (N-terminus) of the polypeptide chain, plays a dual role. It not only targets the ribosome to the ER but also triggers the opening of the translocator. As the protein is fed through the translocator, the signal sequence stays attached, allowing the rest of the protein to move through as a loop. A signal peptidase inside the ER then cuts off the signal sequence, which is subsequently discarded into the lipid bilayer of the ER membrane and broken down. Finally, once the last part of the protein (the C-terminus) passes through the translocator, the entire soluble protein is released into the ER lumen, where it can then fold and undergo further modifications or be transported to its final destination. Mechanisms of transmembrane protein integration Transmembrane proteins, which are partly integrated into the ER membrane rather than released into the ER lumen, have a complex assembly process. The initial stages are similar to soluble proteins: a signal sequence starts the insertion into the ER membrane. However, this process is interrupted by a stop-transfer sequence—a string of hydrophobic amino acids—which causes the translocator to halt and release the protein laterally into the membrane. This results in a single-pass transmembrane protein with one end inside the ER lumen and the other in the cytosol, and this orientation is permanent. Some transmembrane proteins use an internal signal (start-transfer sequence) instead of one at the N-terminus, and unlike the initial signal sequence, this start-transfer sequence isn't removed. It begins the transfer process, which continues until a stop-transfer sequence is encountered, at which point both sequences become anchored in the membrane as alpha-helical segments. In more complex proteins that span the membrane multiple times, additional pairs of start- and stop-transfer sequences are used to weave the protein into the membrane in a fashion akin to a sewing machine. Each pair allows a new segment to cross the membrane and adds to the protein's structure, ensuring it is properly embedded with the correct arrangement of segments inside and outside the ER membrane. Peroxisomes Peroxisomes contain a single phospholipid bilayer that surrounds the peroxisomal matrix containing a wide variety of proteins and enzymes that participate in anabolism and catabolism. Peroxisomes are specialized cell organelles that carry out specific oxidative reactions using molecular oxygen. Their primary function is to remove hydrogen atoms from organic molecules, a process that results in the production of hydrogen peroxide (). Within peroxisomes, an enzyme called catalase plays a critical role. It uses the hydrogen peroxide generated in the earlier reaction to oxidize various other substances, including phenols, formic acid, formaldehyde, and alcohol. This is known as the "peroxidative" reaction. Peroxisomes are particularly important in liver and kidney cells for detoxifying harmful substances that enter the bloodstream. For example, they are responsible for oxidizing about 25% of the ethanol we consume into acetaldehyde. Additionally, catalase within peroxisomes can break down excess hydrogen peroxide into water and oxygen and thus preventing potential damage from the build-up of . Since it contains no internal DNA like that of the mitochondria or chloroplast all peroxisomal proteins are encoded by nuclear genes. To date there are two types of known Peroxisome Targeting Signals (PTS): Peroxisome targeting signal 1 (PTS1): a C-terminal tripeptide with a consensus sequence (S/A/C)-(K/R/H)-(L/A). The most common PTS1 is serine-lysine-leucine (SKL). The initial research that led to the discovery of this consensus observed that when firefly luciferase was expressed in cultured insect cells it was targeted to the peroxisome. By testing a variety of mutations in the gene encoding the expressed luciferase, the consensus sequence was then determined. It has also been found that by adding this C-terminal sequence of SKL to a cytosolic protein that it becomes targeted for transport to the peroxisome. The majority of peroxisomal matrix proteins possess this PTS1 type signal. Peroxisome targeting signal 2 (PTS2): a nonapeptide located near the N-terminus with a consensus sequence (R/K)-(L/V/I)-XXXXX-(H/Q)-(L/A/F) (where X can be any amino acid). There are also proteins that possess neither of these signals. Their transport may be based on a so-called "piggy-back" mechanism: such proteins associate with PTS1-possessing matrix proteins and are translocated into the peroxisomal matrix together with them. In the case of cytosolic proteins that are produced with the PTS1 C-terminal sequence, its path to the peroxisomal matrix is dependent upon binding to another cytosolic protein called pex5 (peroxin 5). Once bound, pex5 interacts with a peroxisomal membrane protein pex14 to form a complex. When the pex5 protein with bound cargo interacts with the pex14 membrane protein, the complex induces the release of the targeted protein into the matrix. Upon releasing the cargo protein into the matrix, pex5 dissociation from pex14 occurs via ubiquitinylation by a membrane complex comprising pex2, pex12, and pex10 followed by an ATP dependent removal involving the cytosolic protein complex pex1 and pex6. The cycle for pex5 mediated import into the peroxisomal matrix is restored after the ATP dependent removal of ubiquitin and is free to bind with another protein containing a PTS1 sequence. Proteins containing a PTS2 targeting sequence are mediated by a different cytosolic protein but are believed to follow a similar mechanism to that of those containing the PTS1 sequence. Diseases Protein transport is defective in the following genetic diseases: Mohr–Tranebjaerg syndrome Zellweger syndrome Adrenoleukodystrophy (ALD). Refsum disease Parkinson's disease Hypercholesterolemia, atherosclerosis, obesity, and diabetes In bacteria and archaea As discussed above (see protein translocation), most prokaryotic membrane-bound and secretory proteins are targeted to the plasma membrane by either a co-translation pathway that uses bacterial SRP or a post-translation pathway that requires SecA and SecB. At the plasma membrane, these two pathways deliver proteins to the SecYEG translocon for translocation. Bacteria may have a single plasma membrane (Gram-positive bacteria), or an inner membrane plus an outer membrane separated by the periplasm (Gram-negative bacteria). Besides the plasma membrane the majority of prokaryotes lack membrane-bound organelles as found in eukaryotes, but they may assemble proteins onto various types of inclusions such as gas vesicles and storage granules. Gram-negative bacteria In gram-negative bacteria proteins may be incorporated into the plasma membrane, the outer membrane, the periplasm or secreted into the environment. Systems for secreting proteins across the bacterial outer membrane may be quite complex and play key roles in pathogenesis. These systems may be described as type I secretion, type II secretion, etc. Gram-positive bacteria In most gram-positive bacteria, certain proteins are targeted for export across the plasma membrane and subsequent covalent attachment to the bacterial cell wall. A specialized enzyme, sortase, cleaves the target protein at a characteristic recognition site near the protein C-terminus, such as an LPXTG motif (where X can be any amino acid), then transfers the protein onto the cell wall. Several analogous systems are found that likewise feature a signature motif on the extra-cytoplasmic face, a C-terminal transmembrane domain, and cluster of basic residues on the cytosolic face at the protein's extreme C-terminus. The PEP-CTERM/exosortase system, found in many Gram-negative bacteria, seems to be related to extracellular polymeric substance production. The PGF-CTERM/archaeosortase A system in archaea is related to S-layer production. The GlyGly-CTERM/rhombosortase system, found in the Shewanella, Vibrio, and a few other genera, seems involved in the release of proteases, nucleases, and other enzymes. Bioinformatic tools Minimotif Miner is a bioinformatics tool that searches protein sequence queries for a known protein targeting sequence motifs. Phobius predicts signal peptides based on a supplied primary sequence. SignalP predicts signal peptide cleavage sites. LOCtree predicts the subcellular localization of proteins.	Biology and health sciences	Cell processes	Biology
24838	https://en.wikipedia.org/wiki/Peptidoglycan	Peptidoglycan	Peptidoglycan or murein is a unique large macromolecule, a polysaccharide, consisting of sugars and amino acids that forms a mesh-like layer (sacculus) that surrounds the bacterial cytoplasmic membrane. The sugar component consists of alternating residues of β-(1,4) linked N-acetylglucosamine (NAG) and N-acetylmuramic acid (NAM). Attached to the N-acetylmuramic acid is an oligopeptide chain made of three to five amino acids. The peptide chain can be cross-linked to the peptide chain of another strand forming the 3D mesh-like layer. Peptidoglycan serves a structural role in the bacterial cell wall, giving structural strength, as well as counteracting the osmotic pressure of the cytoplasm. This repetitive linking results in a dense peptidoglycan layer which is critical for maintaining cell form and withstanding high osmotic pressures, and it is regularly replaced by peptidoglycan production. Peptidoglycan hydrolysis and synthesis are two processes that must occur in order for cells to grow and multiply, a technique carried out in three stages: clipping of current material, insertion of new material, and re-crosslinking of existing material to new material. The peptidoglycan layer is substantially thicker in gram-positive bacteria (20 to 80 nanometers) than in gram-negative bacteria (7 to 8 nanometers). Depending on pH growth conditions, the peptidoglycan forms around 40 to 90% of the cell wall's dry weight of gram-positive bacteria but only around 10% of gram-negative strains. Thus, presence of high levels of peptidoglycan is the primary determinant of the characterisation of bacteria as gram-positive. In gram-positive strains, it is important in attachment roles and serotyping purposes. For both gram-positive and gram-negative bacteria, particles of approximately 2 nm can pass through the peptidoglycan. It is difficult to tell whether an organism is gram-positive or gram-negative using a microscope; Gram staining, created by Hans Christian Gram in 1884, is required. The bacteria are stained with the dyes crystal violet and safranin. Gram positive cells are purple after staining, while Gram negative cells stain pink. Structure The peptidoglycan layer within the bacterial cell wall is a crystal lattice structure formed from linear chains of two alternating amino sugars, namely N-acetylglucosamine (GlcNAc or NAG) and N-acetylmuramic acid (MurNAc or NAM). The alternating sugars are connected by a β-(1,4)-glycosidic bond. Each MurNAc is attached to a short (4- to 5-residue) amino acid chain, containing L-alanine, D-glutamic acid, meso-diaminopimelic acid, and D-alanine in the case of Escherichia coli (a gram-negative bacterium); or L-alanine, D-glutamine, L-lysine, and D-alanine with a 5-glycine interbridge between tetrapeptides in the case of Staphylococcus aureus (a gram-positive bacterium). Peptidoglycan is one of the most important sources of D-amino acids in nature. By enclosing the inner membrane, the peptidoglycan layer protects the cell from lysis caused by the turgor pressure of the cell. When the cell wall grows, it retains its shape throughout its life, so a rod shape will remain a rod shape, and a spherical shape will remain a spherical shape for life. This happens because the freshly added septal material of synthesis transforms into a hemispherical wall for the offspring cells. Cross-linking between amino acids in different linear amino sugar chains occurs with the help of the enzyme DD-transpeptidase and results in a 3-dimensional structure that is strong and rigid. The specific amino acid sequence and molecular structure vary with the bacterial species. The different peptidoglycan types of bacterial cell walls and their taxonomic implications have been described. Archaea (domain Archaea) do not contain peptidoglycan (murein). Some Archaea contain pseudopeptidoglycan (pseudomurein, see below). Peptidoglycan is involved in binary fission during bacterial cell reproduction. L-form bacteria and mycoplasmas, both lacking peptidoglycan cell walls, do not proliferate by binary fission, but by a budding mechanism. In the course of early evolution, the successive development of boundaries (membranes, walls) protecting first structures of life against their environment must have been essential for the formation of the first cells (cellularisation). The invention of rigid peptidoglycan (murein) cell walls in bacteria (domain Bacteria) was probably the prerequisite for their survival, extensive radiation and colonisation of virtually all habitats of the geosphere and hydrosphere. Biosynthesis The peptidoglycan monomers are synthesized in the cytosol and are then attached to a membrane carrier bactoprenol. Bactoprenol transports peptidoglycan monomers across the cell membrane where they are inserted into the existing peptidoglycan. In the first step of peptidoglycan synthesis, glutamine, which is an amino acid, donates an amino group to a sugar, fructose 6-phosphate. This reaction, catalyzed by EC 2.6.1.16 (GlmS), turns fructose 6-phosphate into glucosamine-6-phosphate. In step two, an acetyl group is transferred from acetyl CoA to the amino group on the glucosamine-6-phosphate creating N-acetyl-glucosamine-6-phosphate. This reaction is EC 5.4.2.10, catalyzed by GlmM. In step three of the synthesis process, the N-acetyl-glucosamine-6-phosphate is isomerized, which will change N-acetyl-glucosamine-6-phosphate to N-acetyl-glucosamine-1-phosphate. This is EC 2.3.1.157, catalyzed by GlmU. In step 4, the N-acetyl-glucosamine-1-phosphate, which is now a monophosphate, attacks UTP. Uridine triphosphate, which is a pyrimidine nucleotide, has the ability to act as an energy source. In this particular reaction, after the monophosphate has attacked the UTP, an inorganic pyrophosphate is given off and is replaced by the monophosphate, creating UDP-N-acetylglucosamine (2,4). (When UDP is used as an energy source, it gives off an inorganic phosphate.) This initial stage, is used to create the precursor for the NAG in peptidoglycan. This is EC 2.7.7.23, also catalyzed by GlmU, which is a bifunctional enzyme. In step 5, some of the UDP-N-acetylglucosamine (UDP-GlcNAc) is converted to UDP-MurNAc (UDP-N-acetylmuramic acid) by the addition of a lactyl group to the glucosamine. Also in this reaction, the C3 hydroxyl group will remove a phosphate from the alpha carbon of phosphoenolpyruvate. This creates what is called an enol derivative. EC 2.5.1.7, catalyzed by MurA. In step 6, the enol is reduced to a "lactyl moiety" by NADPH in step six. EC 1.3.1.98, catalyzed by MurB. In step 7, the UDP–MurNAc is converted to UDP-MurNAc pentapeptide by the addition of five amino acids, usually including the dipeptide D-alanyl-D-alanine. This is a string of three reactions: EC 6.3.2.8 by MurC, EC 6.3.2.9 by MurD, and EC 6.3.2.13 by MurE. Each of these reactions requires the energy source ATP. This is all referred to as Stage one. Stage two occurs in the cytoplasmic membrane. It is in the membrane where a lipid carrier called bactoprenol carries peptidoglycan precursors through the cell membrane. Undecaprenyl phosphate will attack the UDP-MurNAc penta, creating a PP-MurNac penta, which is now a lipid (lipid I). EC 2.7.8.13 by MraY. UDP-GlcNAc is then transported to MurNAc, creating Lipid-PP-MurNAc penta-GlcNAc (lipid II), a disaccharide, also a precursor to peptidoglycan. EC 2.4.1.227 by MurG. Lipid II is transported across the membrane by flippase (MurJ), a discovery made in 2014 after decades of searching. Once it is there, it is added to the growing glycan chain by the enzyme peptidoglycan glycosyltransferase (GTase, EC 2.4.1.129). This reaction is known as transglycosylation. In the reaction, the hydroxyl group of the GlcNAc will attach to the MurNAc in the glycan, which will displace the lipid-PP from the glycan chain. In a final step, the DD-transpeptidase (TPase, EC 3.4.16.4) crosslinks individual glycan chains. This protein is also known as the penicillin-binding protein. Some versions of the enzyme also performs the glycosyltransferase function, while others leave the job to a separate enzyme. Pseudopeptidoglycan In some archaea, i.e. members of the Methanobacteriales and in the genus Methanopyrus, pseudopeptidoglycan (pseudomurein) has been found. In pseudopeptidoglycan the sugar residues are β-(1,3) linked N-acetylglucosamine and N-acetyltalosaminuronic acid. This makes the cell walls of such archaea insensitive to lysozyme. The biosynthesis of pseudopeptidoglycan has been described. Recognition by immune system Peptidoglycan recognition is an evolutionarily conserved process. The overall structure is similar between bacterial species, but various modifications can increase the diversity. These include modifications of the length of sugar polymers, modifications in the sugar structures, variations in cross-linking or substitutions of amino acids (primarily at the third position). The aim of these modifications is to alter the properties of the cell wall, which plays a vital role in pathogenesis. Peptidoglycans can be degraded by several enzymes (lysozyme, glucosaminidase, endopeptidase...), producing immunostimulatory fragments (sometimes called muropeptides) that are critical for mediating host-pathogen interactions. These include MDP (muramyl dipeptide), NAG (N-acetylglucosamine) or iE-DAP (γ-d-glutamyl-meso-diaminopimelic acid). Peptidoglycan from intestinal bacteria (both pathogens and commensals) crosses the intestinal barrier even under physiological conditions. Mechanisms through which peptidoglycan or its fragments enter the host cells can be direct (carrier-independent) or indirect (carrier-dependent), and they are either bacteria-mediated (secretion systems, membrane vesicles) or host cell-mediated (receptor-mediated, peptide transporters). Bacterial secretion systems are protein complexes used for the delivery of virulence factors across the bacterial cell envelope to the exterior environment. Intracellular bacterial pathogens invade eukaryotic cells (which may lead to the formation of phagolysosomes and/or autophagy activation), or bacteria may be engulfed by phagocytes (macrophages, monocytes, neutrophils...). The bacteria-containing phagosome may then fuse with endosomes and lysosomes, leading to degradation of bacteria and generation of polymeric peptidoglycan fragments and muropeptides. Receptors Innate immune system senses intact peptidoglycan and peptidoglycan fragments using numerous PRRs (pattern recognition receptors) that are secreted, expressed intracellularly or expressed on the cell surface. Peptidoglycan recognition proteins PGLYRPs are conserved from insects to mammals. Mammals produce four secreted soluble peptidoglycan recognition proteins (PGLYRP-1, PGLYRP-2, PGLYRP-3 and PGLYRP-4) that recognize muramyl pentapeptide or tetrapeptide. They can also bind to LPS and other molecules by using binding sites outside of the peptidoglycan-binding groove. After recognition of peptidoglycan, PGLYRPs activate polyphenol oxidase (PPO) molecules, Toll, or immune deficiency (IMD) signalling pathways. That leads to production of antimicrobial peptides (AMPs). Each of the mammalian PGLYRPs display unique tissue expression patterns. PGLYRP-1 is mainly expressed in the granules of neutrophils and eosinophils. PGLYRP-3 and 4 are expressed by several tissues such as skin, sweat glands, eyes or the intestinal tract. PGLYRP-1, 3 and 4 form disulphide-linked homodimers and heterodimers essential for their bactericidal activity. Their binding to bacterial cell wall peptidoglycans can induce bacterial cell death by interaction with various bacterial transcriptional regulatory proteins. PGLYRPs are likely to assist in bacterial killing by cooperating with other PRRs to enhance recognition of bacteria by phagocytes. PGLYRP-2 is primarily expressed by the liver and secreted into the circulation. Also, its expression can be induced in skin keratinocytes, oral and intestinal epithelial cells. In contrast with the other PGLYRPs, PGLYRP-2 has no direct bactericidal activity. It possesses peptidoglycan amidase activity, it hydrolyses the lactyl-amide bond between the MurNAc and the first amino acid of the stem peptide of peptidoglycan. It is proposed, that the function of PGLYRP-2 is to prevent over-activation of the immune system and inflammation-induced tissue damage in response to NOD2 ligands (see below), as these muropeptides can no longer be recognized by NOD2 upon separation of the peptide component from MurNAc. Growing evidence suggests that peptidoglycan recognition protein family members play a dominant role in the tolerance of intestinal epithelial cells toward the commensal microbiota. It has been demonstrated that expression of PGLYRP-2 and 4 can influence the composition of the intestinal microbiota. Recently, it has been discovered, that PGLYRPs (and also NOD-like receptors and peptidoglycan transporters) are highly expressed in the developing mouse brain. PGLYRP-2 and is highly expressed in neurons of several brain regions including the prefrontal cortex, hippocampus, and cerebellum, thus indicating potential direct effects of peptidoglycan on neurons. PGLYRP-2 is highly expressed also in the cerebral cortex of young children, but not in most adult cortical tissues. PGLYRP-1 is also expressed in the brain and continues to be expressed into adulthood. NOD-like receptors Probably the most well-known receptors of peptidoglycan are the NOD-like receptors (NLRs), mainly NOD1 and NOD2. The NOD1 receptor is activated after iE-DAP (γ-d-glutamyl-meso-diaminopimelic acid) binding, while NOD2 recognizes MDP (muramyl dipeptide), by their LRR domains. Activation leads to self-oligomerization, resulting in activation of two signalling cascades. One triggers activation of NF-κB (through RIP2, TAK1 and IKK), second leads to MAPK signalling cascade. Activation of these pathways induces production of inflammatory cytokines and chemokines. NOD1 is expressed by diverse cell types, including myeloid phagocytes, epithelial cells and neurons. NOD2 is expressed in monocytes and macrophages, epithelial intestinal cells, Paneth cells, dendritic cells, osteoblasts, keratinocytes and other epithelial cell types. As cytosolic sensors, NOD1 and NOD2 must either detect bacteria that enter the cytosol, or peptidoglycan must be degraded to generate fragments that must be transported into the cytosol for these sensors to function. Recently, it was demonstrated that NLRP3 is activated by peptidoglycan, through a mechanism that is independent of NOD1 and NOD2. In macrophages, N-acetylglucosamine generated by peptidoglycan degradation was found to inhibit hexokinase activity and induce its release from the mitochondrial membrane. It promotes NLRP3 inflammasome activation through a mechanism triggered by increased mitochondrial membrane permeability. NLRP1 is also considered as a cytoplasmic sensor of peptidoglycan. It can sense MDP and promote IL-1 secretion through binding NOD2. C-type lectin receptors (CLRs) C-type lectins are a diverse superfamily of mainly Ca2+-dependent proteins that bind a variety of carbohydrates (including the glycan skeleton of peptidoglycan), and function as innate immune receptors. CLR proteins that bind to peptidoglycan include MBL (mannose binding lectin), ficolins, Reg3A (regeneration gene family protein 3A) and PTCLec1. In mammals, they initiate the lectin-pathway of the complement cascade. Toll-like receptors The role of TLRs in direct recognition of peptidoglycan is controversial. In some studies, has been reported that peptidoglycan is sensed by TLR2. But this TLR2-inducing activity could be due to cell wall lipoproteins and lipoteichoic acids that commonly co-purify with peptidoglycan. Also variation in peptidoglycan structure in bacteria from species to species may contribute to the differing results on this topic. As vaccine or adjuvant Peptidoglycan is immunologically active, which can stimulate immune cells to increase the expression of cytokines and enhance antibody-dependent specific response when combined with vaccine or as adjuvant alone. MDP, which is the basic unit of peptidoglycan, was initially used as the active component of Freund's adjuvant. Peptidoglycan from Staphylococcus aureus was used as a vaccine to protect mice, showing that after vaccine injection for 40 weeks, the mice survived from S. aureus challenge at an increased lethal dose. Inhibition and degradation Some antibacterial drugs such as penicillin interfere with the production of peptidoglycan by binding to bacterial enzymes known as penicillin-binding proteins or DD-transpeptidases. Penicillin-binding proteins form the bonds between oligopeptide crosslinks in peptidoglycan. For a bacterial cell to reproduce through binary fission, more than a million peptidoglycan subunits (NAM-NAG+oligopeptide) must be attached to existing subunits. Mutations in genes coding for transpeptidases that lead to reduced interactions with an antibiotic are a significant source of emerging antibiotic resistance. Since peptidoglycan is also lacking in L-form bacteria and in mycoplasmas, both are resistant against penicillin. Other steps of peptidoglycan synthesis can also be targeted. The topical antibiotic bacitracin targets the utilization of C55-isoprenyl pyrophosphate. Lantibiotics, which includes the food preservative nisin, attack lipid II. Lysozyme, which is found in tears and constitutes part of the body's innate immune system exerts its antibacterial effect by breaking the β-(1,4)-glycosidic bonds in peptidoglycan (see above). Lysozyme is more effective in acting against gram-positive bacteria, in which the peptidoglycan cell wall is exposed, than against gram-negative bacteria, which have an outer layer of LPS covering the peptidoglycan layer. Several bacterial peptidoglycan modifications can result in resistance to degradation by lysozyme. Susceptibility of bacteria to degradation is also considerably affected by exposure to antibiotics. Exposed bacteria synthesize peptidoglycan that contains shorter sugar chains that are poorly crosslinked and this peptidoglycan is then more easily degraded by lysozyme.	Biology and health sciences	Proteins	Biology
24863	https://en.wikipedia.org/wiki/Pseudomonadota	Pseudomonadota	Pseudomonadota (synonym Proteobacteria) is a major phylum of Gram-negative bacteria. Currently, they are considered the predominant phylum within the realm of bacteria. They are naturally found as pathogenic and free-living (non-parasitic) genera. The phylum comprises six classes Acidithiobacillia, Alphaproteobacteria, Betaproteobacteria, Gammaproteobacteria, Hydrogenophilia, and Zetaproteobacteria. The Pseudomonadota are widely diverse, with differences in morphology, metabolic processes, relevance to humans, and ecological influence. Classification American microbiologist Carl Woese established this grouping in 1987, calling it informally the "purple bacteria and their relatives". The group was later formally named the 'Proteobacteria' after the Greek god Proteus, who was known to assume many forms. In 2021 the International Committee on Systematics of Prokaryotes designated the synonym Pseudomonadota, and renamed many other prokaryotic phyla as well. This renaming of several prokaryote phyla in 2021, including Pseudomonadota, remains controversial among microbiologists, many of whom continue to use the earlier name Proteobacteria, of long standing in the literature. The phylum Pseudomonadota encompasses classes Acidithiobacillia, Alphaproteobacteria, Betaproteobacteria, Gammaproteobacteria, Hydrogenophilia, and Zetaproteobacteria. The phylum includes a wide variety of pathogenic genera, such as Escherichia, Salmonella, Vibrio, Yersinia, Legionella, and many others. Others are free-living (non-parasitic) and include many of the bacteria responsible for nitrogen fixation. Previously, the Pseudomonadota phylum included two additional classes, namely Deltaproteobacteria and Oligoflexia. However, further investigation into the phylogeny of these taxa through genomic marker analysis demonstrated their separation from the Pseudomonadota phylum. Deltaproteobacteria has been identified as a diverse taxonomic unit, leading to a proposal for its reclassification into distinct phyla: Desulfobacterota (encompassing Thermodesulfobacteria), Myxococcota, and Bdellovibrionota (comprising Oligoflexia). The class Epsilonproteobacteria was additionally identified within the Pseudomonadota phylum. This class is characterized by its significance as chemolithotrophic primary producers and its metabolic prowess in deep-sea hydrothermal vent ecosystems. Noteworthy pathogenic genera within this class include Campylobacter, Helicobacter, and Arcobacter. Analysis of phylogenetic tree topology and genetic markers revealed the direct divergence of Epsilonproteobacteria from the Pseudomonadota phylum. Limited outgroup data and low bootstrap values support these discoveries. Despite further investigations, consensus has not been reached regarding the monophyletic nature of Epsilonproteobacteria within Proteobacteria, prompting researchers to propose its taxonomic separation from the phylum. The proposed reclassification of the name Epsilonproteobacteria is Epsilonbacteraeota, later revised to Campylobacterota in 2018. Taxonomy The currently accepted taxonomy is based on the List of Prokaryotic names with Standing in Nomenclature (LSPN) and the National Center for Biotechnology Information (NCBI). The group Pseudomonadota is defined based on ribosomal RNA (rRNA) sequencing, and are divided into several subclasses. These subclasses were regarded as such for many years, but are now treated as various classes of the phylum. These classes are monophyletic. The genus Acidithiobacillus, part of the Gammaproteobacteria until it was transferred to class Acidithiobacillia in 2013, was previously regarded as paraphyletic to the Betaproteobacteria according to multigenome alignment studies. In 2017, the Betaproteobacteria was subject to major revisions and the class Hydrogenophilalia was created to contain the order Hydrogenophilales Pseudomonadota classes with validly published names include some prominent genera: e.g.: Acidithiobacillia: Acidithiobacillus, Thermithiobacillus Alphaproteobacteria: Brucella, Rhizobium, Agrobacterium, Caulobacter, Rickettsia, Wolbachia, etc. Betaproteobacteria: Bordetella, Ralstonia, Neisseria, Nitrosomonas, etc. Gammaproteobacteria: Escherichia, Shigella, Salmonella, Yersinia, Buchnera, Haemophilus, Vibrio, Pseudomonas, Pasteurella, etc. Zetaproteobacteria: Mariprofundus Characteristics Pseudomonadota are a diverse group. Though some species may stain Gram-positive or Gram-variable in the laboratory, they are nominally Gram-negative. Their unique outer membrane is mainly composed of lipopolysaccharides, which helps differentiate them from the Gram-positive species. Most Pseudomonadota are motile and move using flagella. Many move about using flagella, but some are nonmotile, or rely on bacterial gliding. Pseudomonadota have a wide variety of metabolism types. Most are facultative or obligate anaerobes, chemolithoautotrophs, and heterotrophs, though numerous exceptions exist. A variety of distantly related genera within the Pseudomonadota obtain their energy from light through conventional photosynthesis or anoxygenic photosynthesis. The Acidithiobacillia contain only sulfur, iron, and uranium-oxidizing autotrophs. The type order is the Acidithiobacillaceae, which includes five different Acidithiobacillus species used in the mining industry. In particular, these microbes assist with the process of bioleaching, which involves microbes assisting in metal extraction from mining waste that typically extraction methods cannot remove. Some Alphaproteobacteria can grow at very low levels of nutrients and have unusual morphology within their life cycles. Some form stalks to help with colonization, and form buds during cell division. Others include agriculturally important bacteria capable of inducing nitrogen fixation in symbiosis with plants. The type order is the Caulobacterales, comprising stalk-forming bacteria such as Caulobacter. The mitochondria of eukaryotes are thought to be descendants of an alphaproteobacterium. The Betaproteobacteria are highly metabolically diverse and contain chemolithoautotrophs, photoautotrophs, and generalist heterotrophs. The type order is the Burkholderiales, comprising an enormous range of metabolic diversity, including opportunistic pathogens. These pathogens are primary for both humans and animals, such as the horse pathogen Burkholderia mallei, and Burkholderia cepacia which causes respiratory tract infections in people with cystic fibrosis. The Gammaproteobacteria are one of the largest classes in terms of genera, containing approximately 250 validly published names. The type order is the Pseudomonadales, which include the genera Pseudomonas and the nitrogen-fixing Azotobacter, along with many others. Besides being a well-known pathogenic genus, Pseudomonas is also capable of biodegradation of certain materials, like cellulose. The Hydrogenophilalia are thermophilic chemoheterotrophs and autotrophs. The bacteria typically use hydrogen gas as an electron donor, but can also use reduced sulfuric compounds. Because of this ability, scientists have begun to use certain species of Hydrogenophilalia to remove sulfides that contaminate industrial wastewater systems. The type order is the Hydrogenophilaceae which contains the genera Thiobacillus, Petrobacter, Sulfuricella, Hydrogenophilus and Tepidiphilus. Currently, no members of this class have been identified as pathogenic. The Zetaproteobacteria are the iron-oxidizing neutrophilic chemolithoautotrophs, distributed worldwide in estuaries and marine habitats. This group is so successful in its environment due to their microaerophilic nature. Because they require less oxygen than what is present in the atmosphere, they are able to compete with the abiotic iron(II) oxidation that is already occurring in the environment. The only confirmed type order for this class is the Mariprofundaceae, which does not contain any known pathogenic species. Transformation Transformation, a process in which genetic material passes from one bacterium to another, has been reported in at least 30 species of Pseudomonadota distributed in the classes alpha, beta, and gamma. The best-studied Pseudomonadota with respect to natural genetic transformation are the medically important human pathogens Neisseria gonorrhoeae (class beta), and Haemophilus influenzae (class gamma). Natural genetic transformation is a sexual process involving DNA transfer from one bacterial cell to another through the intervening medium and the integration of the donor sequence into the recipient genome. In pathogenic Pseudomonadota, transformation appears to serve as a DNA repair process that protects the pathogen's DNA from attack by their host's phagocytic defenses that employ oxidative free radicals. Habitat Due to the distinctive nature of each of the six classes of Pseudomonadota, this phylum occupies a multitude of habitats. These include: Human oral cavity Microbial mats in the deep sea Marine sediments Thermal sulfur springs Agricultural soil Hydrothermal vents Stem nodules of legumes Within aphids as endosymbionts Gastrointestinal tract of warm-blooded species Brackish, estuary waters Microbiomes of shrimp and mollusks Human vaginal tract Potato rhizosphere microbiome Significance Human health Studies have suggested Pseudomonadota as a relevant signature of disease in the human gastrointestinal (GI) tract, by operating as a marker for microbiota instability. The human gut microbiome consists mainly of four phyla: Firmicutes, Bacteroidetes, Actinobacteria, and Pseudomonadota. Microorganism gut colonization is dynamic from birth to death, with stabilization at the first few years of life, to higher diversity in adults, to reduced diversity in the elderly. The gut microbiome conducts processes like nutrient synthesis, chemical metabolism, and the formation of the gut barrier. Additionally, the gut microbiome facilitates host interactions with its surrounding environment through regulation of nutrient absorption and bacterial intake. In 16s rRNA and metagenome sequencing studies, Proteobacteria have been identified as bacteria that prompts endotoxemia (an inflammatory gut response) and metabolic disorders in human GI tracts. Another study by Michail et al. showed a correlation of microbial composition in children with and without nonalcoholic fatty liver disease (NAFLD), wherein patients with NAFLD have a higher abundance of Gammaproteobacteria than patients without the disease. Classes Betaproteobacteria and Gammaproteobacteria are prevalent within the human oral cavity, and are markers for good oral health. The oral microbiome consists of 11 habitats, including the tongue dorsum, hard palate, tonsils, throat, saliva, and more. Changes in the oral microbiome are due to endogenous and exogenous factors like host lifestyle, genotype, environment, immune system, and socioeconomic status. Considering diet as a factor, high saturated fatty acid (SAF) content, achieved through poor diet, has been correlated to increased abundance of Betaproteobacteria in the oral cavity. Economic value Pseudomonadota bacteria have a symbiotic or mutualistic association with plant roots, an example being in the rhizomes of potato plants. Because of this symbiotic relationship, farmers have the ability to increase their crop yields. Healthier root systems can lead to better nutrient uptake, improved water retention, increased resistance to diseases and pests, and ultimately higher crop yields per acre. Increased agricultural output can spark economic growth, contribute to food security, and lead to job creation in rural areas. As briefly mentioned in previous sections, members of Pseudomonadota have vast metabolic abilities that allow them to utilize and produce a variety of compounds. Bioleaching, done by various Thiobacillus species, are a primary example of this. Any iron and sulfur oxidizing species has the potential to uncover metals and low-grade ores that conventional mining techniques were unable to extract. At present, they are most often used for recovering copper and uranium, but researchers are looking to expand this field in the future. The downside of this method is that the bacteria produce acidic byproducts that end up in acid mine drainage. Bioleaching has significant economic promise if it can be controlled and not cause any further harm to the environment. Ecological impact Pseudomonadota are microbes commonly found within soil systems. Microbes play a crucial role in the surrounding ecosystem by performing functions such as nutrient cycling, carbon dioxide fixation, decomposition, and nitrogen fixation. Pseudomonadota can be described as phototrophs, heterotrophs, and lithotrophs. As heterotrophs (examples Pseudomonas and Xanthomonas) these bacteria are effective in breaking down organic matter, contributing to nutrient cycling. Additionally, photolithotrophs within the phylum are able to perform photosynthesis using sulfide or elemental sulfur as electron donors, which enables them to participate in carbon fixation and oxygen production even in anaerobic conditions. These Pseudomonadota bacteria are also considered copiotrophic organisms, meaning they can be found in environments with high nutrient availability. These environments have ample sources of carbon and other nutrients, environments like fertile soils, compost, and sewage. These copiotrophic bacteria are able to enhance soil health by performing nutrient cycling and waste decomposition. Because this phylum are able to form a symbiotic relationship with plant roots, incorporating Pseudomonadota into agricultural practices aligns with principles of sustainable farming. These bacteria contribute to soil health and fertility, promote natural pest management, and enhance the resilience of crops to environmental stressors.	Biology and health sciences	Gram-negative bacteria	Plants
24868	https://en.wikipedia.org/wiki/Pauli%20matrices	Pauli matrices	In mathematical physics and mathematics, the Pauli matrices are a set of three complex matrices that are traceless, Hermitian, involutory and unitary. Usually indicated by the Greek letter sigma (), they are occasionally denoted by tau () when used in connection with isospin symmetries. These matrices are named after the physicist Wolfgang Pauli. In quantum mechanics, they occur in the Pauli equation, which takes into account the interaction of the spin of a particle with an external electromagnetic field. They also represent the interaction states of two polarization filters for horizontal/vertical polarization, 45 degree polarization (right/left), and circular polarization (right/left). Each Pauli matrix is Hermitian, and together with the identity matrix (sometimes considered as the zeroth Pauli matrix ), the Pauli matrices form a basis of the vector space of Hermitian matrices over the real numbers, under addition. This means that any Hermitian matrix can be written in a unique way as a linear combination of Pauli matrices, with all coefficients being real numbers. The Pauli matrices satisfy the useful product relation: Hermitian operators represent observables in quantum mechanics, so the Pauli matrices span the space of observables of the complex two-dimensional Hilbert space. In the context of Pauli's work, represents the observable corresponding to spin along the th coordinate axis in three-dimensional Euclidean space The Pauli matrices (after multiplication by to make them anti-Hermitian) also generate transformations in the sense of Lie algebras: the matrices form a basis for the real Lie algebra , which exponentiates to the special unitary group SU(2). The algebra generated by the three matrices is isomorphic to the Clifford algebra of and the (unital) associative algebra generated by functions identically (is isomorphic) to that of quaternions (). Algebraic properties All three of the Pauli matrices can be compacted into a single expression: where the solution to is the "imaginary unit", and is the Kronecker delta, which equals if and 0 otherwise. This expression is useful for "selecting" any one of the matrices numerically by substituting values of in turn useful when any of the matrices (but no particular one) is to be used in algebraic manipulations. The matrices are involutory: where is the identity matrix. The determinants and traces of the Pauli matrices are from which we can deduce that each matrix has eigenvalues +1 and −1. With the inclusion of the identity matrix (sometimes denoted ), the Pauli matrices form an orthogonal basis (in the sense of Hilbert–Schmidt) of the Hilbert space of Hermitian matrices over , and the Hilbert space of all complex matrices over . Commutation and anti-commutation relations Commutation relations The Pauli matrices obey the following commutation relations: where the Levi-Civita symbol is used. These commutation relations make the Pauli matrices the generators of a representation of the Lie algebra Anticommutation relations They also satisfy the anticommutation relations: where is defined as and is the Kronecker delta. denotes the identity matrix. These anti-commutation relations make the Pauli matrices the generators of a representation of the Clifford algebra for denoted The usual construction of generators of using the Clifford algebra recovers the commutation relations above, up to unimportant numerical factors. A few explicit commutators and anti-commutators are given below as examples: Eigenvectors and eigenvalues Each of the (Hermitian) Pauli matrices has two eigenvalues: and . The corresponding normalized eigenvectors are Pauli vectors The Pauli vector is defined by where , , and are an equivalent notation for the more familiar , , and . The Pauli vector provides a mapping mechanism from a vector basis to a Pauli matrix basis as follows: More formally, this defines a map from to the vector space of traceless Hermitian matrices. This map encodes structures of as a normed vector space and as a Lie algebra (with the cross-product as its Lie bracket) via functions of matrices, making the map an isomorphism of Lie algebras. This makes the Pauli matrices intertwiners from the point of view of representation theory. Another way to view the Pauli vector is as a Hermitian traceless matrix-valued dual vector, that is, an element of that maps Completeness relation Each component of can be recovered from the matrix (see completeness relation below) This constitutes an inverse to the map , making it manifest that the map is a bijection. Determinant The norm is given by the determinant (up to a minus sign) Then, considering the conjugation action of an matrix on this space of matrices, we find and that is Hermitian and traceless. It then makes sense to define where has the same norm as and therefore interpret as a rotation of three-dimensional space. In fact, it turns out that the special restriction on implies that the rotation is orientation preserving. This allows the definition of a map given by where This map is the concrete realization of the double cover of by and therefore shows that The components of can be recovered using the tracing process above: Cross-product The cross-product is given by the matrix commutator (up to a factor of ) In fact, the existence of a norm follows from the fact that is a Lie algebra (see Killing form). This cross-product can be used to prove the orientation-preserving property of the map above. Eigenvalues and eigenvectors The eigenvalues of are This follows immediately from tracelessness and explicitly computing the determinant. More abstractly, without computing the determinant, which requires explicit properties of the Pauli matrices, this follows from since this can be factorised into A standard result in linear algebra (a linear map that satisfies a polynomial equation written in distinct linear factors is diagonal) means this implies is diagonal with possible eigenvalues The tracelessness of means it has exactly one of each eigenvalue. Its normalized eigenvectors are These expressions become singular for . They can be rescued by letting and taking the limit , which yields the correct eigenvectors (0,1) and (1,0) of . Alternatively, one may use spherical coordinates to obtain the eigenvectors and . Pauli 4-vector The Pauli 4-vector, used in spinor theory, is written with components given in terms of the ???, and the ?????, , This defines a map from to the vector space of Hermitian matrices, which also encodes the Minkowski metric (with mostly minus convention) in its determinant: This 4-vector also has a completeness relation. It is convenient to define a second Pauli 4-vector and allow raising and lowering using the Minkowski metric tensor. The relation can then be written Similarly to the Pauli 3-vector case, we can find a matrix group that acts as isometries on in this case the matrix group is and this shows Similarly to above, this can be explicitly realized for with components In fact, the determinant property follows abstractly from trace properties of the For matrices, the following identity holds: That is, the 'cross-terms' can be written as traces. When are chosen to be different the cross-terms vanish. It then follows, now showing summation explicitly, Since the matrices are this is equal to Relation to dot and cross product Pauli vectors elegantly map these commutation and anticommutation relations to corresponding vector products. Adding the commutator to the anticommutator gives so that, Contracting each side of the equation with components of two -vectors and (which commute with the Pauli matrices, i.e., for each matrix and vector component (and likewise with ) yields Finally, translating the index notation for the dot product and cross product results in If is identified with the pseudoscalar then the right hand side becomes , which is also the definition for the product of two vectors in geometric algebra. If we define the spin operator as , then satisfies the commutation relation:Or equivalently, the Pauli vector satisfies: Some trace relations The following traces can be derived using the commutation and anticommutation relations. If the matrix is also considered, these relationships become where Greek indices and assume values from and the notation is used to denote the sum over the cyclic permutation of the included indices. Exponential of a Pauli vector For one has, for even powers, which can be shown first for the case using the anticommutation relations. For convenience, the case is taken to be by convention. For odd powers, Matrix exponentiating, and using the Taylor series for sine and cosine, . In the last line, the first sum is the cosine, while the second sum is the sine; so, finally, which is analogous to Euler's formula, extended to quaternions. Note that , while the determinant of the exponential itself is just , which makes it the generic group element of SU(2). A more abstract version of formula for a general matrix can be found in the article on matrix exponentials. A general version of for an analytic (at a and −a) function is provided by application of Sylvester's formula, The group composition law of A straightforward application of formula provides a parameterization of the composition law of the group . One may directly solve for in which specifies the generic group multiplication, where, manifestly, the spherical law of cosines. Given , then, Consequently, the composite rotation parameters in this group element (a closed form of the respective BCH expansion in this case) simply amount to (Of course, when is parallel to , so is , and .) Adjoint action It is also straightforward to likewise work out the adjoint action on the Pauli vector, namely rotation of any angle along any axis : Taking the dot product of any unit vector with the above formula generates the expression of any single qubit operator under any rotation. For example, it can be shown that . Completeness relation An alternative notation that is commonly used for the Pauli matrices is to write the vector index in the superscript, and the matrix indices as subscripts, so that the element in row and column of the -th Pauli matrix is In this notation, the completeness relation for the Pauli matrices can be written As noted above, it is common to denote the 2 × 2 unit matrix by so The completeness relation can alternatively be expressed as The fact that any Hermitian complex 2 × 2 matrices can be expressed in terms of the identity matrix and the Pauli matrices also leads to the Bloch sphere representation of 2 × 2 mixed states’ density matrix, (positive semidefinite 2 × 2 matrices with unit trace. This can be seen by first expressing an arbitrary Hermitian matrix as a real linear combination of as above, and then imposing the positive-semidefinite and trace conditions. For a pure state, in polar coordinates, the idempotent density matrix acts on the state eigenvector with eigenvalue +1, hence it acts like a projection operator. Relation with the permutation operator Let be the transposition (also known as a permutation) between two spins and living in the tensor product space This operator can also be written more explicitly as Dirac's spin exchange operator, Its eigenvalues are therefore 1 or −1. It may thus be utilized as an interaction term in a Hamiltonian, splitting the energy eigenvalues of its symmetric versus antisymmetric eigenstates. SU(2) The group SU(2) is the Lie group of unitary matrices with unit determinant; its Lie algebra is the set of all anti-Hermitian matrices with trace 0. Direct calculation, as above, shows that the Lie algebra is the three-dimensional real algebra spanned by the set . In compact notation, As a result, each can be seen as an infinitesimal generator of SU(2). The elements of SU(2) are exponentials of linear combinations of these three generators, and multiply as indicated above in discussing the Pauli vector. Although this suffices to generate SU(2), it is not a proper representation of , as the Pauli eigenvalues are scaled unconventionally. The conventional normalization is so that As SU(2) is a compact group, its Cartan decomposition is trivial. SO(3) The Lie algebra is isomorphic to the Lie algebra , which corresponds to the Lie group SO(3), the group of rotations in three-dimensional space. In other words, one can say that the are a realization (and, in fact, the lowest-dimensional realization) of infinitesimal rotations in three-dimensional space. However, even though and are isomorphic as Lie algebras, and are not isomorphic as Lie groups. is actually a double cover of , meaning that there is a two-to-one group homomorphism from to , see relationship between SO(3) and SU(2). Quaternions The real linear span of is isomorphic to the real algebra of quaternions, , represented by the span of the basis vectors The isomorphism from to this set is given by the following map (notice the reversed signs for the Pauli matrices): Alternatively, the isomorphism can be achieved by a map using the Pauli matrices in reversed order, As the set of versors forms a group isomorphic to , gives yet another way of describing . The two-to-one homomorphism from to may be given in terms of the Pauli matrices in this formulation. Physics Classical mechanics In classical mechanics, Pauli matrices are useful in the context of the Cayley-Klein parameters. The matrix corresponding to the position of a point in space is defined in terms of the above Pauli vector matrix, Consequently, the transformation matrix for rotations about the -axis through an angle may be written in terms of Pauli matrices and the unit matrix as Similar expressions follow for general Pauli vector rotations as detailed above. Quantum mechanics In quantum mechanics, each Pauli matrix is related to an angular momentum operator that corresponds to an observable describing the spin of a spin particle, in each of the three spatial directions. As an immediate consequence of the Cartan decomposition mentioned above, are the generators of a projective representation (spin representation) of the rotation group SO(3) acting on non-relativistic particles with spin . The states of the particles are represented as two-component spinors. In the same way, the Pauli matrices are related to the isospin operator. An interesting property of spin particles is that they must be rotated by an angle of 4 in order to return to their original configuration. This is due to the two-to-one correspondence between SU(2) and SO(3) mentioned above, and the fact that, although one visualizes spin up/down as the north–south pole on the 2-sphere they are actually represented by orthogonal vectors in the two-dimensional complex Hilbert space. For a spin particle, the spin operator is given by , the fundamental representation of SU(2). By taking Kronecker products of this representation with itself repeatedly, one may construct all higher irreducible representations. That is, the resulting spin operators for higher spin systems in three spatial dimensions, for arbitrarily large j, can be calculated using this spin operator and ladder operators. They can be found in . The analog formula to the above generalization of Euler's formula for Pauli matrices, the group element in terms of spin matrices, is tractable, but less simple. Also useful in the quantum mechanics of multiparticle systems, the general Pauli group is defined to consist of all -fold tensor products of Pauli matrices. Relativistic quantum mechanics In relativistic quantum mechanics, the spinors in four dimensions are 4 × 1 (or 1 × 4) matrices. Hence the Pauli matrices or the Sigma matrices operating on these spinors have to be 4 × 4 matrices. They are defined in terms of 2 × 2 Pauli matrices as It follows from this definition that the matrices have the same algebraic properties as the matrices. However, relativistic angular momentum is not a three-vector, but a second order four-tensor. Hence needs to be replaced by , the generator of Lorentz transformations on spinors. By the antisymmetry of angular momentum, the are also antisymmetric. Hence there are only six independent matrices. The first three are the The remaining three, where the Dirac matrices are defined as The relativistic spin matrices are written in compact form in terms of commutator of gamma matrices as Quantum information In quantum information, single-qubit quantum gates are 2 × 2 unitary matrices. The Pauli matrices are some of the most important single-qubit operations. In that context, the Cartan decomposition given above is called the "Z–Y decomposition of a single-qubit gate". Choosing a different Cartan pair gives a similar "X–Y decomposition of a single-qubit gate.	Physical sciences	Quantum mechanics	Physics
24873	https://en.wikipedia.org/wiki/Polearm	Polearm	A polearm or pole weapon is a close combat weapon in which the main fighting part of the weapon is fitted to the end of a long shaft, typically of wood, extending the user's effective range and striking power. Polearms are predominantly melee weapons, with a subclass of spear-like designs fit for thrusting and/or throwing. Because many polearms were adapted from agricultural implements or other fairly abundant tools, and contained relatively little metal, they were cheap to make and readily available. When belligerents in warfare had a poorer class who could not pay for dedicated military weapons, they would often appropriate tools as cheap weapons. The cost of training was comparatively low, since these conscripted farmers had spent most of their lives using these "weapons" in the fields. This made polearms the favoured weapon of peasant levies and peasant rebellions the world over. Polearms can be divided into three broad categories: those designed for extended reach and thrusting tactics used in pike square or phalanx combat; those designed to increase leverage (due to hands moving freely on a pole) to maximize angular force (swinging tactics) against cavalry; and those designed for throwing tactics used in skirmish line combat. The hook on weapons such as the halberd was used for pulling or grappling tactics, especially against horsemen. Because of their versatility, high effectiveness and low cost, there were many variants of polearm, which were much-used weapons on the battlefield. Bills, picks, dane axes, spears, glaives, guandaos, pudaos, pikes, poleaxes, halberds, harpoons, sovnyas, tridents, naginatas, bardiches, war scythes, and lances are all varieties of polearms. Polearms were common weapons on post-classical battlefields of Asia and Europe. Their range and impact force made them effective weapons against armoured warriors on horseback, unhorsing the opponent and to some extent effective to penetrate armour. The Renaissance saw a plethora of varieties. Polearms in modern times are largely constrained to ceremonial military units such as the Papal Swiss Guard or Yeomen of the Guard, or traditional martial arts. Chinese martial arts in particular have preserved a wide variety of weapons and techniques. Classification difficulties The classification of polearms can be difficult, and European weapon classifications in particular can be confusing. This can be due to a number of factors, including uncertainty in original descriptions, changes in weapons or nomenclature through time, mistranslation of terms, and the well-meaning inventiveness of later experts. For example, the word "halberd" is also used to translate the Chinese ji and also a range of medieval Scandinavian weapons as described in sagas, such as the atgeir. As well, all polearms developed from three early tools (the axe, the scythe, and the knife) and one weapon, the spear. In the words of the arms expert Ewart Oakeshott, While men-at-arms may have been armed with custom designed military weapons, militias were often armed with whatever was available. These may or may not have been mounted on poles and described by one of more names. The problems with precise definitions can be inferred by a contemporary description of Royalist infantry which were engaged in the Battle of Birmingham (1643) during the first year of English Civil War (in the early modern period). The infantry regiment that accompanied Prince Rupert's cavalry were armed: List of polearms Ancient polearms European Dory (spear) Falx Kontos (weapon) Rhomphaia Sarissa Trident Xyston Asian Dagger-axe The dagger-axe (Chinese: 戈; pinyin: gē; Wade–Giles: ko; sometimes confusingly translated "halberd") is a type of weapon that was in use from Shang dynasty until at least Han dynasty China. It consists of a dagger-shaped blade made of bronze (or later iron) mounted by the tang to a perpendicular wooden shaft: a common Bronze Age infantry weapon, also used by charioteers. Some dagger axes include a spear-point. There is a (rare) variant type with a divided two-part head, consisting of the usual straight blade and a scythe-like blade. Other rarities include archaeology findings with two or sometimes three blades stacked in line on top of a pole, but were generally thought as ceremonial polearms. Though the weapon saw frequent use in ancient China, the use of the dagger-axe decreased dramatically after the Qin and Han dynasties. The ji combines the dagger axe with a spear. By the post-classical Chinese dynasties, with the decline of chariot warfare, the use of the dagger-axe was almost nonexistent. Ji The ji (Chinese: 戟) was created by combining the dagger-axe with a spear. It was used as a military weapon at least as early as the Shang dynasty until the end of the Northern and Southern dynasties. Ngao The ngao or ngau (ง้าว,ของ้าว) is a Thai polearm that was traditionally used by elephant-riding infantry and is still used by practitioners of krabi krabong. Known in Malay as a dap, it consists of a wooden shaft with a curved blade fashioned onto the end, and is similar in design to the Korean woldo. Usually, it also had a hook (ขอ) between the blade and shaft used for commanding the elephant. The elephant warrior used the ngao like a blade from atop an elephant or horse during battle. Post-classical polearms European Dane axe The Dane axe is a weapon with a heavy crescent-shaped head mounted on a haft in length. Originally a Viking weapon, it was adopted by the Anglo-Saxons and Normans in the 11th century, spreading through Europe in the 12th and 13th centuries. Variants of this basic weapon continued in use in Scotland and Ireland into the 16th century. A form of 'long axe'. Sparth axe In the 13th century, variants on the Danish axe are seen. Described in English as a "sparth" (from the Old Norse ) or "pale-axe", the weapon featured a larger head with broader blade, the rearward part of the crescent sweeping up to contact (or even be attached to) the haft. In Ireland, this axe was known as a "sparr axe". Originating in either Western Scotland or Ireland, the sparr was widely used by the galloglass. Although sometimes said to derive from the Irish for a joist or beam, a more likely definition is as a variant of sparth. Although attempts have been made to suggest that the sparr had a distinctive shaped head, illustrations and surviving weapons show there was considerable variation and the distinctive feature of the weapon was its long haft. Fauchard A fauchard is a type of polearm which was used in medieval Europe from the 11th through the 14th centuries. The design consists of a curved blade put atop a pole. The blade bears a moderate to strong curve along its length; however, unlike a bill or guisarme, the cutting edge is on the convex side. Guisarme A guisarme (sometimes gisarme, giserne or bisarme) is a polearm used in Europe primarily between 1000 and 1400. It was used primarily to dismount knights and horsemen. Like most polearms it was developed by peasants by combining hand tools with long poles, in this case by putting a pruning hook onto a spear shaft. While early designs were simply a hook on the end of a long pole, later designs implemented a small reverse spike on the back of the blade. Eventually weapon makers incorporated the usefulness of the hook in a variety of different polearms and guisarme became a catch-all for any weapon that included a hook on the blade. Ewart Oakeshott has proposed an alternative description of the weapon as a crescent shaped socketed axe. Glaive A glaive is a polearm consisting of a single-edged tapering blade similar in shape to a modern kitchen knife on the end of a pole. The blade was around long, on the end of a pole long. However, instead of having a tang like a sword or naginata, the blade is affixed in a socket-shaft configuration similar to an axe head, both the blade and shaft varying in length. Illustrations in the 13th century Maciejowski Bible show a short staffed weapon with a long blade used by both infantry and cavalry. Occasionally glaive blades were created with a small hook or spike on the reverse side. Such glaives are named glaive-guisarme. Voulge A voulge (occasionally called a pole cleaver) is a curved blade attached to a pole by binding the lower two-thirds of the blade to the side of the pole, to form a sort of axe. Looks very similar to a glaive. Svärdstav A svärdstav (literally sword-staff) is a Swedish medieval polearm that consists of a two-edged sword blade attached to a staff. The illustrations often show the weapon being equipped with sword-like quillons. The illustrations sometimes show a socket mount and reinforcing langets being used, but sometimes they are missing; it is possible this weapon was sometimes manufactured by simply attaching an old sword blade onto a long pole on its tang, not unlike a naginata. Asian Naginata A naginata (なぎなた or 薙刀) is a Japanese polearm that was traditionally used by members of the samurai class. A naginata consists of a wood shaft with a curved blade on the end. Usually it also had a sword-like guard (tsuba) between the blade and shaft. It was mounted with a tang and held in place with a pin or pins, rather than going over the shaft using a socket. The naginata was developed based on the hoko yari from the 1st millennium AD or the tachi from the late Heian period (794ー1185). It was appreciated by samurai who fought on foot as a weapon to maintain optimal distance from the enemy in close combat, but after the Onin War in the 15th century, large groups of mobilized infantry called asigaru began to equip themselves with yari (spear) yumi (longbow) and tanegashima (gun), making naginata and tachi (long sword) obsolete on the battlefield and often replaced with nagamaki and katana. From the Edo period, naginata has been recognized as a martial art practiced by women in the samurai class. Yari A yari (やり or 槍) is a Japanese polearm that was traditionally used by members of the samurai class. There are various types of yari, which have different names depending on the shape of the blade attached to the end of the wooden shaft. For example, 'Jumonji yari' refers to a yari with a cross-shaped blade, and 'Sasaho yari' refers to a yari with a blade shaped like a sasa leaf. During the Sengoku period, a large group of ashigaru in a formation used yari as one of their main weapons and exerted tremendous power on the battlefield. Honda Tadakatsu a vassal of Tokugawa Ieyasu, had gained a reputation as a master of one of the Three Great Spears of Japan, Tonbokiri. Woldo The Korean woldo was a variation of the Chinese guan dao. It was originally used by the post-classical Shilla warriors. Wielding the woldo took time due to its weight, but in the hands of a trained soldier, the woldo was a fearsome, agile weapon famous for enabling a single soldier to cut down ranks of infantrymen. The woldo was continually in use for the military in Korea with various modifications made over the decades. Unlike the Chinese with the guan dao, the Koreans found the woldo unwieldy on horseback, and thus, it was specifically tailored to the needs of infantrymen. The Joseon government implemented rigorous training regimens requiring soldiers to be proficient with swordsmanship, and the use of the woldo. Though it was never widely used as a standard weapon, the woldo saw action on many fronts and was considered by many Korean troops to be a versatile weapon. Recently, a contemporary revival in various martial arts in Korea has brought interest into the application of the woldo and its history. Guandao A guandao or kwan tou is a type of Chinese polearm. In Chinese, it is properly called a yanyue dao (偃月刀), 'reclining moon blade'. Some believed it comes from the late Han Era and was supposedly used by the late Eastern Han dynasty general Guan Yu, but archaeological findings have shown that Han dynasty armies generally used straight, single-edged blades, and curved blades came several centuries later. There is no reason to believe their polearms had curved blades on them. Besides, historical accounts of the Three Kingdoms era describe Guan Yu thrusting his opponents down (probably with a spear-like polearm) in battle, not cutting them down with a curved blade. The guandao is also known as the chun qiu da dao ('spring autumn great knife'), again probably related to the depiction of Guan Yu in the Ming dynasty novel Romance of the Three Kingdoms, but possibly a Ming author's invention. It consists of a heavy blade mounted atop a wooden or metal pole with a pointed metal counter weight used for striking and stabbing on the opposite end. The blade is very deep and curved on its face, resembling a Chinese saber, or dao. Variant designs include rings along the length of the straight back edge, as found in the nine-ring guandao. The "elephant" guandao's tip curls into a rounded spiral, while the dragon head guandao features a more ornate design. Podao A podao, 'long-handled sabre', is a Chinese polearm, also known as the zhan ma dao ('horsecutter sabre'), which has a lighter blade and a ring at the end. A podao is an infantryman's weapon, mainly used for cutting the legs off oncoming charging horses to bring down the riders. Fangtian ji In the Song dynasty, several weapons were referred to as ji, but they were developed from spears, not from ancient ji. One variety was called the qinglong ji (), and had a spear tip with a crescent blade on one side. Another type was the fangtian ji (), which had a spear tip with crescent blades on both sides. They had multiple means of attack: the side blade or blades, the spear tip, plus often a rear counterweight that could be used to strike the opponent. The way the side blades were fixed to the shaft differs, but usually there were empty spaces between the pole and the side blade. The wielder could strike with the shaft, with the option of then pulling the weapon back to hook with a side blade; or, he could slap his opponent with the flat side of the blade to knock him off his horse. Barcha and Ballam The Barcha is a type of lance with a wooden handle, once common in South Asia in the 16th century and was popular weapon of choice in the Maratha Empire. Variations of the barcha is the hand-like Karpa Barcha and the serpent-like Nagni Barcha. Another variant included the Ballam, a javelin effective at bringing down infantry and cavalry at a distance. Nagni Barcha is identified as the weapon used by the Sikh warrior Bhai Bachittar Singh to kill a drunken Mughal war elephant at the Siege of Lohgarh. Later polearms European Bardiche Bec de Corbin Bill Bohemian earspoon Brandistock Lochaber axe Lucerne hammer Military fork Partisan Pike Ranseur Scottish polearms Sovnya Spetum Viking halberd War scythe Corseque A corseque has a three-bladed head on a haft which, like the partisan, is similar to the winged spear or spetum in the later Middle Ages. It was popular in Europe in the 16th and 17th centuries. Surviving examples have a variety of head forms but there are two main variants, one with the side blades (known as flukes or wings) branching from the neck of the central blade at 45 degrees, the other with hooked blades curving back towards the haft. The corseque is usually associated with the rawcon, ranseur and runka. Another possible association is with the "three-grayned staff" listed as being in the armoury of Henry VIII in 1547 (though the same list also features 84 rawcons, suggesting the weapons were not identical in 16th century English eyes). Another modern term used for particularly ornate-bladed corseques is the chauve-souris. Halberd A halberd (or Swiss voulge) is a two-handed polearm that came to prominent use during the 14th and 15th centuries but has continued in use as a ceremonial weapon to the present day. First recorded as "hellembart" in 1279, the word halberd possibly comes from the German words Halm (staff) or Helm (helmet), and Barte (axe). The halberd consists of an axe blade topped with a spike mounted on a long shaft. It always has a hook or thorn on the back side of the axe blade for grappling mounted combatants. Early forms are very similar in many ways to certain forms of voulge, while 16th century and later forms are similar to the pollaxe. The Swiss were famous users of the halberd in the medieval and renaissance eras, with various cantons evolving regional variations of the basic form. Poleaxe In the 14th century, the basic long axe gained an armour-piercing spike on the back and another on the end of the haft for thrusting. This is similar to the pollaxe of 15th century. The poleaxe emerged in response to the need for a weapon that could penetrate plate armour and featured various combinations of an axe-blade, a back-spike and a hammer. It was the favoured weapon for men-at-arms fighting on foot into the sixteenth century.	Technology	Polearms	null
24884	https://en.wikipedia.org/wiki/Peppered%20moth	Peppered moth	The peppered moth (Biston betularia) is a temperate species of night-flying moth. It is mostly found in the northern hemisphere in places like Asia, Europe and North America. Peppered moth evolution is an example of population genetics and natural selection. The caterpillars of the peppered moth not only mimic the form but also the colour of a twig. Recent research indicates that the caterpillars can sense the twig's colour with their skin and match their body colour to the background to protect themselves from predators. Description The wingspan ranges from 45 mm to 62 mm (median 55 mm). It is relatively stout-bodied, with forewings relatively narrow-elongate. The wings are white, "peppered" with black, and with more-or-less distinct cross lines, also black. These transverse wing lines and "peppered" maculation (spotting) can also, in rare instances, be gray or brown; the spotting pattern, in particularly very rare cases, is sometimes a combination of brown and black/gray. The black speckling varies in amount, in some examples it is almost absent, whilst in others it is so dense that the wings appear to be black sprinkled with white. The antennae of males are strongly bipectinate. Prout (1912–16) gives an account of the forms and congeners. Distribution Biston betularia is found in China (Heilongjiang, Jilin, Inner Mongolia, Beijing, Hebei, Shanxi, Shandong, Henan, Shaanxi, Ningxia, Gansu, Qinghai, Xinjiang, Fujian, Sichuan, Yunnan, Tibet), Russia, Mongolia, Japan, North Korea, South Korea, Nepal, Kazakhstan, Kyrgyzstan, Turkmenistan, Georgia, Azerbaijan, Armenia, Europe and North America. Ecology and life cycle In Great Britain and Ireland, the peppered moth is univoltine (i.e., it has one generation per year), whilst in south-eastern North America it is bivoltine (two generations per year). The lepidopteran life cycle consists of four stages: ova (eggs), several larval instars (caterpillars), pupae, which overwinter in the soil, and imagines (adults). During the day, the moths typically rest on trees, where they are preyed on by birds. The caterpillar is a twig mimic, varying in colour between green and brown. On a historical note, it was one of the first animals to be identified as being camouflaged with countershading to make it appear flat (shading being the main visual cue that makes things appear solid), in a paper by Edward Bagnall Poulton in 1887. Research indicates that the caterpillars can sense the twig's colour with their skin and match their body colour to the background to protect themselves from predators, an ability to camouflage themselves also found in cephalopods, chameleons and some fish, although this colour change is rather slower in the caterpillars. It goes into the soil late in the season, where it pupates in order to spend the winter. The imagines emerge from the pupae between late May and August, the males slightly before the females (this is common and expected from sexual selection). They emerge late in the day and dry their wings before flying that night. The males fly every night of their lives in search of females, whereas the females only fly on the first night. Thereafter, the females release pheromones to attract males. Since the pheromone is carried by the wind, males tend to travel up the concentration gradient, i.e., toward the source. During flight, they are subject to predation by bats. The males guard the female from other males until she lays the eggs. The female lays about 2,000 pale-green ovoid eggs about 1 mm in length into crevices in bark with her ovipositor. Resting behaviour A mating pair or a lone individual will spend the day hiding from predators, particularly birds. In the case of the former, the male stays with the female to ensure paternity. Evidence for resting positions is given by data collected by the peppered moth researcher Michael Majerus, and it is given in the accompanying charts. These data were originally published in Howlett and Majerus (1987), and an updated version published in Majerus (1998), who concluded that the moths rest in the upper part of the trees. Majerus notes: Creationist critics of the peppered moth have often pointed to a statement made by Clarke et al. (1985): "... In 25 years we have only found two betularia on the tree trunks or walls adjacent to our traps, and none elsewhere". The reason now seems obvious. Few people spend their time looking for moths up in the trees. That is where peppered moths rest by day. From their original data, Howlett and Majerus (1987) concluded that peppered moths generally rest in unexposed positions, using three main types of site. Firstly, a few inches below a branch-trunk joint on a tree trunk where the moth is in shadow; secondly, on the underside of branches and thirdly on foliate twigs. The above data would appear to support this. Further support for these resting positions is given from experiments watching captive moths taking up resting positions in both males (Mikkola, 1979; 1984) and females (Liebert and Brakefield, 1987). Majerus, et al., (2000) have shown that peppered moths are cryptically camouflaged against their backgrounds when they rest in the boughs of trees. It is clear that in human visible wavelengths, typica are camouflaged against lichens and carbonaria against plain bark. However, birds are capable of seeing ultraviolet light that humans cannot see. Using an ultraviolet-sensitive video camera, Majerus et al. showed that typica reflect ultraviolet light in a speckled fashion and are camouflaged against crustose lichens common on branches, both in ultraviolet and human-visible wavelengths. However, typica are not as well camouflaged against foliose lichens common on tree trunks; though they are camouflaged in human wavelengths, in ultraviolet wavelengths, foliose lichens do not reflect ultraviolet light. During an experiment in Cambridge over the seven years 2001–2007 Majerus noted the natural resting positions of peppered moths, and of the 135 moths examined over half were on tree branches, mostly on the lower half of the branch, 37% were on tree trunks, mostly on the north side, and only 12.6% were resting on or under twigs. Polymorphism Introduction on forms There are several melanic and non-melanic morphs of the peppered moth. These are controlled genetically. A particular colour morph can be indicated in a standard way by following the species name in the form "morpha morph name". The use of "form" in the method of Biston betularia f. formname in detailing these variations is also a widespread practice. These forms are often accidentally elevated to subspecies status when they appear in literature. Not adding the "f." (forma) or morpha implies that the taxon is a subspecies instead of a form, as in Biston betularia carbonaria instead of Biston betularia f. carbonaria. Rarely, forms have been elevated to species status, as in Biston carbonaria. Either of these two circumstances might lead to the erroneous belief that speciation was involved in the observed evolution of the peppered moth. This is not the case; individuals of each morph interbreed and produce fertile offspring with individuals of all other morphs; hence there is only one peppered moth species. By contrast, different subspecies of the same species can theoretically interbreed with one another and will produce fully fertile and healthy offspring, but in practice do not, as they live in different regions or reproduce in different seasons. Full-fledged species are either unable to produce fertile and healthy offspring, or do not recognize each other's courtship signals, or both. European breeding experiments have shown that in Biston betularia betularia, the allele for melanism producing morpha carbonaria is controlled by a single locus. The melanic allele is dominant to the non-melanic allele. This situation is, however, somewhat complicated by the presence of three other alleles that produce indistinguishable morphs of morpha medionigra. These are of intermediate dominance, but this is not complete (Majerus, 1998). Form names In continental Europe, there are three morphs: the white morph typica (syn. morpha/f. betularia), the dark melanistic morph carbonaria (syn. doubledayaria), and an intermediate form medionigra. In Britain, the typical white morph is known as typica, the melanic morph is carbonaria, and the intermediate phenotype is named insularia. In North America, the melanic black morph is morpha swettaria. In Biston betularia cognataria, the melanic allele (producing morpha swettaria) is similarly dominant to the non-melanic allele. There are also some intermediate morphs. In Japan, no melanic morphs have been recorded; they are all morpha typica. Evolution The evolution of the peppered moth over the last two hundred years has been studied in detail. At the start of this period, the vast majority of peppered moths had light coloured wing patterns which effectively camouflaged them against the light-coloured trees and lichens upon which they rested. However, due to widespread pollution during the Industrial Revolution in England, many of the lichens died out, and the trees which peppered moths rested on became blackened by soot, causing most of the light-coloured moths, or typica, to die off due to predation. At the same time, the dark-coloured, or melanic, moths, carbonaria, flourished because they could hide on the darkened trees. Since then, with improved environmental standards, light-coloured peppered moths have again become common, and the dramatic change in the peppered moth's population has remained a subject of much interest and study. This has led to the coining of the term "industrial melanism" to refer to the genetic darkening of species in response to pollutants. As a result of the relatively simple and easy-to-understand circumstances of the adaptation, the peppered moth has become a common example used in explaining or demonstrating natural selection to laypeople and classroom students through simulations. The first carbonaria morph was recorded by Edleston in Manchester in 1848, and over the subsequent years it increased in frequency. Predation experiments, particularly by Bernard Kettlewell, established that the agent of selection was birds who preyed on the carbonaria morph. Subsequent experiments and observations have supported the initial evolutionary explanation of the phenomenon. Genetic basis of melanism The evolution of the industrial melanism mutation has been shown to be due to the insertion of a transposable element into the first intron of the cortex gene, resulting in an increase in the abundance of the cortex transcript, which is expressed in developing wings. Gallery	Biology and health sciences	Lepidoptera	Animals
24888	https://en.wikipedia.org/wiki/Promoter%20%28genetics%29	Promoter (genetics)	In genetics, a promoter is a sequence of DNA to which proteins bind to initiate transcription of a single RNA transcript from the DNA downstream of the promoter. The RNA transcript may encode a protein (mRNA), or can have a function in and of itself, such as tRNA or rRNA. Promoters are located near the transcription start sites of genes, upstream on the DNA (towards the 5' region of the sense strand). Promoters can be about 100–1000 base pairs long, the sequence of which is highly dependent on the gene and product of transcription, type or class of RNA polymerase recruited to the site, and species of organism. Overview For transcription to take place, the enzyme that synthesizes RNA, known as RNA polymerase, must attach to the DNA near a gene. Promoters contain specific DNA sequences such as response elements that provide a secure initial binding site for RNA polymerase and for proteins called transcription factors that recruit RNA polymerase. These transcription factors have specific activator or repressor sequences of corresponding nucleotides that attach to specific promoters and regulate gene expression. In bacteria The promoter is recognized by RNA polymerase and an associated sigma factor, which in turn are often brought to the promoter DNA by an activator protein's binding to its own DNA binding site nearby. In eukaryotes The process is more complicated, and at least seven different factors are necessary for the binding of an RNA polymerase II to the promoter. Promoters represent critical elements that can work in concert with other regulatory regions (enhancers, silencers, boundary elements/insulators) to direct the level of transcription of a given gene. A promoter is induced in response to changes in abundance or conformation of regulatory proteins in a cell, which enable activating transcription factors to recruit RNA polymerase. Given the short sequences of most promoter elements, promoters can rapidly evolve from random sequences. For instance, in E. coli, ~60% of random sequences can evolve expression levels comparable to the wild-type lac promoter with only one mutation, and that ~10% of random sequences can serve as active promoters even without evolution. Identification of relative location As promoters are typically immediately adjacent to the gene in question, positions in the promoter are designated relative to the transcriptional start site, where transcription of DNA begins for a particular gene (i.e., positions upstream are negative numbers counting back from -1, for example -100 is a position 100 base pairs upstream). Elements Bacterial In bacteria, the promoter contains two short sequence elements approximately 10 (Pribnow Box) and 35 nucleotides upstream from the transcription start site. The sequence at -10 (the -10 element) has the consensus sequence TATAAT. The sequence at -35 (the -35 element) has the consensus sequence TTGACA. The above consensus sequences, while conserved on average, are not found intact in most promoters. On average, only 3 to 4 of the 6 base pairs in each consensus sequence are found in any given promoter. Few natural promoters have been identified to date that possess intact consensus sequences at both the -10 and -35; artificial promoters with complete conservation of the -10 and -35 elements have been found to transcribe at lower frequencies than those with a few mismatches with the consensus. The optimal spacing between the -35 and -10 sequences is 17 bp. Some promoters contain one or more upstream promoter element (UP element) subsites (consensus sequence 5'-AAAAAARNR-3' when centered in the -42 region; consensus sequence 5'-AWWWWWTTTTT-3' when centered in the -52 region; W = A or T; R = A or G; N = any base). The above promoter sequences are recognized only by RNA polymerase holoenzyme containing sigma-70. RNA polymerase holoenzymes containing other sigma factors recognize different core promoter sequences. ← upstream downstream → 5'-XXXXXXXPPPPPPXXXXXXPPPPPPXXXXGGGGGGGGGGGGGGGGGGGGGGGGGGGGGGGGGGGGGGGGGGGGGGGGGGGGGXXXX-3' -35 -10 Gene to be transcribed Probability of occurrence of each nucleotide for -10 sequence T A T A A T 77% 76% 60% 61% 56% 82% for -35 sequence T T G A C A 69% 79% 61% 56% 54% 54% Bidirectional (prokaryotic) Promoters can be very closely located in the DNA. Such "closely spaced promoters" have been observed in the DNAs of all life forms, from humans to prokaryotes and are highly conserved. Therefore, they may provide some (presently unknown) advantages. These pairs of promoters can be positioned in divergent, tandem, and convergent directions. They can also be regulated by transcription factors and differ in various features, such as the nucleotide distance between them, the two promoter strengths, etc. The most important aspect of two closely spaced promoters is that they will, most likely, interfere with each other. Several studies have explored this using both analytical and stochastic models. There are also studies that measured gene expression in synthetic genes or from one to a few genes controlled by bidirectional promoters. More recently, one study measured most genes controlled by tandem promoters in E. coli. In that study, two main forms of interference were measured. One is when an RNAP is on the downstream promoter, blocking the movement of RNAPs elongating from the upstream promoter. The other is when the two promoters are so close that when an RNAP sits on one of the promoters, it blocks any other RNAP from reaching the other promoter. These events are possible because the RNAP occupies several nucleotides when bound to the DNA, including in transcription start sites. Similar events occur when the promoters are in divergent and convergent formations. The possible events also depend on the distance between them. Eukaryotic Gene promoters are typically located upstream of the gene and can have regulatory elements several kilobases away from the transcriptional start site (enhancers). In eukaryotes, the transcriptional complex can cause the DNA to bend back on itself, which allows for placement of regulatory sequences far from the actual site of transcription. Eukaryotic RNA-polymerase-II-dependent promoters can contain a TATA box (consensus sequence TATAAA), which is recognized by the general transcription factor TATA-binding protein (TBP); and a B recognition element (BRE), which is recognized by the general transcription factor TFIIB. The TATA element and BRE typically are located close to the transcriptional start site (typically within 30 to 40 base pairs). Eukaryotic promoter regulatory sequences typically bind proteins called transcription factors that are involved in the formation of the transcriptional complex. An example is the E-box (sequence CACGTG), which binds transcription factors in the basic helix-loop-helix (bHLH) family (e.g. BMAL1-Clock, cMyc). Some promoters that are targeted by multiple transcription factors might achieve a hyperactive state, leading to increased transcriptional activity. Core promoter – the minimal portion of the promoter required to properly initiate transcription Includes the transcription start site (TSS) and elements directly upstream A binding site for RNA polymerase RNA polymerase I: transcribes genes encoding 18S, 5.8S and 28S ribosomal RNAs RNA polymerase II: transcribes genes encoding messenger RNA and certain small nuclear RNAs and microRNA RNA polymerase III: transcribes genes encoding transfer RNA, 5s ribosomal RNAs and other small RNAs General transcription factor binding sites, e.g. TATA box, B recognition element. Many other elements/motifs may be present. There is no such thing as a set of "universal elements" found in every core promoter. Proximal promoter – the proximal sequence upstream of the gene that tends to contain primary regulatory elements Approximately 250 base pairs upstream of the start site Specific transcription factor binding sites Distal promoter – the distal sequence upstream of the gene that may contain additional regulatory elements, often with a weaker influence than the proximal promoter Anything further upstream (but not an enhancer or other regulatory region whose influence is positional/orientation independent) Specific transcription factor binding sites Mammalian promoters Up-regulated expression of genes in mammals is initiated when signals are transmitted to the promoters associated with the genes. Promoter DNA sequences may include different elements such as CpG islands (present in about 70% of promoters), a TATA box (present in about 24% of promoters), initiator (Inr) (present in about 49% of promoters), upstream and downstream TFIIB recognition elements (BREu and BREd) (present in about 22% of promoters), and downstream core promoter element (DPE) (present in about 12% of promoters). The presence of multiple methylated CpG sites in CpG islands of promoters causes stable silencing of genes. However, the presence or absence of the other elements have relatively small effects on gene expression in experiments. Two sequences, the TATA box and Inr, caused small but significant increases in expression (45% and 28% increases, respectively). The BREu and the BREd elements significantly decreased expression by 35% and 20%, respectively, and the DPE element had no detected effect on expression. Cis-regulatory modules that are localized in DNA regions distant from the promoters of genes can have very large effects on gene expression, with some genes undergoing up to 100-fold increased expression due to such a cis-regulatory module. These cis-regulatory modules include enhancers, silencers, insulators and tethering elements. Among this constellation of elements, enhancers and their associated transcription factors have a leading role in the regulation of gene expression. Enhancers are regions of the genome that are major gene-regulatory elements. Enhancers control cell-type-specific gene expression programs, most often by looping through long distances to come in physical proximity with the promoters of their target genes. In a study of brain cortical neurons, 24,937 loops were found, bringing enhancers to promoters. Multiple enhancers, each often at tens or hundred of thousands of nucleotides distant from their target genes, loop to their target gene promoters and coordinate with each other to control expression of their common target gene. The schematic illustration in this section shows an enhancer looping around to come into close physical proximity with the promoter of a target gene. The loop is stabilized by a dimer of a connector protein (e.g. dimer of CTCF or YY1), with one member of the dimer anchored to its binding motif on the enhancer and the other member anchored to its binding motif on the promoter (represented by the red zigzags in the illustration). Several cell function specific transcription factors (there are about 1,600 transcription factors in a human cell) generally bind to specific motifs on an enhancer and a small combination of these enhancer-bound transcription factors, when brought close to a promoter by a DNA loop, govern the level of transcription of the target gene. Mediator (coactivator) (a complex usually consisting of about 26 proteins in an interacting structure) communicates regulatory signals from enhancer DNA-bound transcription factors directly to the RNA polymerase II (pol II) enzyme bound to the promoter. Enhancers, when active, are generally transcribed from both strands of DNA with RNA polymerases acting in two different directions, producing two eRNAs as illustrated in the Figure. An inactive enhancer may be bound by an inactive transcription factor. Phosphorylation of the transcription factor may activate it and that activated transcription factor may then activate the enhancer to which it is bound (see small red star representing phosphorylation of transcription factor bound to enhancer in the illustration). An activated enhancer begins transcription of its RNA before activating a promoter to initiate transcription of messenger RNA from its target gene. Bidirectional (mammalian) Bidirectional promoters are short (<1 kbp) intergenic regions of DNA between the 5' ends of the genes in a bidirectional gene pair. A "bidirectional gene pair" refers to two adjacent genes coded on opposite strands, with their 5' ends oriented toward one another. The two genes are often functionally related, and modification of their shared promoter region allows them to be co-regulated and thus co-expressed. Bidirectional promoters are a common feature of mammalian genomes. About 11% of human genes are bidirectionally paired. Bidirectionally paired genes in the Gene Ontology database shared at least one database-assigned functional category with their partners 47% of the time. Microarray analysis has shown bidirectionally paired genes to be co-expressed to a higher degree than random genes or neighboring unidirectional genes. Although co-expression does not necessarily indicate co-regulation, methylation of bidirectional promoter regions has been shown to downregulate both genes, and demethylation to upregulate both genes. There are exceptions to this, however. In some cases (about 11%), only one gene of a bidirectional pair is expressed. In these cases, the promoter is implicated in suppression of the non-expressed gene. The mechanism behind this could be competition for the same polymerases, or chromatin modification. Divergent transcription could shift nucleosomes to upregulate transcription of one gene, or remove bound transcription factors to downregulate transcription of one gene. Some functional classes of genes are more likely to be bidirectionally paired than others. Genes implicated in DNA repair are five times more likely to be regulated by bidirectional promoters than by unidirectional promoters. Chaperone proteins are three times more likely, and mitochondrial genes are more than twice as likely. Many basic housekeeping and cellular metabolic genes are regulated by bidirectional promoters. The overrepresentation of bidirectionally paired DNA repair genes associates these promoters with cancer. Forty-five percent of human somatic oncogenes seem to be regulated by bidirectional promoters – significantly more than non-cancer causing genes. Hypermethylation of the promoters between gene pairs WNT9A/CD558500, CTDSPL/BC040563, and KCNK15/BF195580 has been associated with tumors. Certain sequence characteristics have been observed in bidirectional promoters, including a lack of TATA boxes, an abundance of CpG islands, and a symmetry around the midpoint of dominant Cs and As on one side and Gs and Ts on the other. A motif with the consensus sequence of TCTCGCGAGA, also called the CGCG element, was recently shown to drive PolII-driven bidirectional transcription in CpG islands. CCAAT boxes are common, as they are in many promoters that lack TATA boxes. In addition, the motifs NRF-1, GABPA, YY1, and ACTACAnnTCCC are represented in bidirectional promoters at significantly higher rates than in unidirectional promoters. The absence of TATA boxes in bidirectional promoters suggests that TATA boxes play a role in determining the directionality of promoters, but counterexamples of bidirectional promoters do possess TATA boxes and unidirectional promoters without them indicates that they cannot be the only factor. Although the term "bidirectional promoter" refers specifically to promoter regions of mRNA-encoding genes, luciferase assays have shown that over half of human genes do not have a strong directional bias. Research suggests that non-coding RNAs are frequently associated with the promoter regions of mRNA-encoding genes. It has been hypothesized that the recruitment and initiation of RNA polymerase II usually begins bidirectionally, but divergent transcription is halted at a checkpoint later during elongation. Possible mechanisms behind this regulation include sequences in the promoter region, chromatin modification, and the spatial orientation of the DNA. Subgenomic A subgenomic promoter is a promoter added to a virus for a specific heterologous gene, resulting in the formation of mRNA for that gene alone. Many positive-sense RNA viruses produce these subgenomic mRNAs (sgRNA) as one of the common infection techniques used by these viruses and generally transcribe late viral genes. Subgenomic promoters range from 24 nucleotide (Sindbis virus) to over 100 nucleotides (Beet necrotic yellow vein virus) and are usually found upstream of the transcription start. Detection A wide variety of algorithms have been developed to facilitate detection of promoters in genomic sequence, and promoter prediction is a common element of many gene prediction methods. A promoter region is located before the -35 and -10 Consensus sequences. The closer the promoter region is to the consensus sequences the more often transcription of that gene will take place. There is not a set pattern for promoter regions as there are for consensus sequences. Binding The initiation of the transcription is a multistep sequential process that involves several mechanisms: promoter location, initial reversible binding of RNA polymerase, conformational changes in RNA polymerase, conformational changes in DNA, binding of nucleoside triphosphate (NTP) to the functional RNA polymerase-promoter complex, and nonproductive and productive initiation of RNA synthesis. The promoter binding process is crucial in the understanding of the process of gene expression. Tuning synthetic genetic systems relies on precisely engineered synthetic promoters with known levels of transcription rates. Location Although RNA polymerase holoenzyme shows high affinity to non-specific sites of the DNA, this characteristic does not allow us to clarify the process of promoter location. This process of promoter location has been attributed to the structure of the holoenzyme to DNA and sigma 4 to DNA complexes. Diseases associated with aberrant function Most diseases are heterogeneous in cause, meaning that one "disease" is often many different diseases at the molecular level, though symptoms exhibited and response to treatment may be identical. How diseases of different molecular origin respond to treatments is partially addressed in the discipline of pharmacogenomics. Not listed here are the many kinds of cancers involving aberrant transcriptional regulation owing to creation of chimeric genes through pathological chromosomal translocation. Importantly, intervention in the number or structure of promoter-bound proteins is one key to treating a disease without affecting expression of unrelated genes sharing elements with the target gene. Some genes whose change is not desirable are capable of influencing the potential of a cell to become cancerous. CpG islands in promoters In humans, about 70% of promoters located near the transcription start site of a gene (proximal promoters) contain a CpG island. CpG islands are generally 200 to 2000 base pairs long, have a C:G base pair content >50%, and have regions of DNA where a cytosine nucleotide is followed by a guanine nucleotide and this occurs frequently in the linear sequence of bases along its 5' → 3' direction. Distal promoters also frequently contain CpG islands, such as the promoter of the DNA repair gene ERCC1, where the CpG island-containing promoter is located about 5,400 nucleotides upstream of the coding region of the ERCC1 gene. CpG islands also occur frequently in promoters for functional noncoding RNAs such as microRNAs. Methylation of CpG islands stably silences genes In humans, DNA methylation occurs at the 5' position of the pyrimidine ring of the cytosine residues within CpG sites to form 5-methylcytosines. The presence of multiple methylated CpG sites in CpG islands of promoters causes stable silencing of genes. Silencing of a gene may be initiated by other mechanisms, but this is often followed by methylation of CpG sites in the promoter CpG island to cause the stable silencing of the gene. Promoter CpG hyper/hypo-methylation in cancer Generally, in progression to cancer, hundreds of genes are silenced or activated. Although silencing of some genes in cancers occurs by mutation, a large proportion of carcinogenic gene silencing is a result of altered DNA methylation (see DNA methylation in cancer). DNA methylation causing silencing in cancer typically occurs at multiple CpG sites in the CpG islands that are present in the promoters of protein coding genes. Altered expressions of microRNAs also silence or activate many genes in progression to cancer (see microRNAs in cancer). Altered microRNA expression occurs through hyper/hypo-methylation of CpG sites in CpG islands in promoters controlling transcription of the microRNAs. Silencing of DNA repair genes through methylation of CpG islands in their promoters appears to be especially important in progression to cancer (see methylation of DNA repair genes in cancer). Canonical sequences and wild-type The usage of the term canonical sequence to refer to a promoter is often problematic, and can lead to misunderstandings about promoter sequences. Canonical implies perfect, in some sense. In the case of a transcription factor binding site, there may be a single sequence that binds the protein most strongly under specified cellular conditions. This might be called canonical. However, natural selection may favor less energetic binding as a way of regulating transcriptional output. In this case, we may call the most common sequence in a population the wild-type sequence. It may not even be the most advantageous sequence to have under prevailing conditions. Recent evidence also indicates that several genes (including the proto-oncogene c-myc) have G-quadruplex motifs as potential regulatory signals. Synthetic promoter design and engineering Promoters are important gene regulatory elements used in tuning synthetically designed genetic circuits and metabolic networks. For example, to overexpress an important gene in a network, to yield higher production of target protein, synthetic biologists design promoters to upregulate its expression. Automated algorithms can be used to design neutral DNA or insulators that do not trigger gene expression of downstream sequences. Diseases that may be associated with variations Some cases of many genetic diseases are associated with variations in promoters or transcription factors. Examples include: Asthma Beta thalassemia Rubinstein-Taybi syndrome Constitutive vs regulated Some promoters are called constitutive as they are active in all circumstances in the cell, while others are regulated, becoming active in the cell only in response to specific stimuli. Tissue-specific promoter A tissue-specific promoter is a promoter that has activity in only certain cell types. Use of the term When referring to a promoter some authors actually mean promoter + operator; i.e., the lac promoter is IPTG inducible, meaning that besides the lac promoter, the lac operon is also present. If the lac operator were not present the IPTG would not have an inducible effect. Another example is the Tac-Promoter system (Ptac). Notice how tac is written as a tac promoter, while in fact tac is actually both a promoter and an operator.	Biology and health sciences	Molecular biology	Biology
24893	https://en.wikipedia.org/wiki/Adobe%20Photoshop	Adobe Photoshop	Adobe Photoshop is a raster graphics editor developed and published by Adobe for Windows and macOS. It was created in 1987 by Thomas and John Knoll. It is the most used tool for professional digital art, especially in raster graphics editing, and its name has become genericised as a verb (e.g. "to photoshop an image", "photoshopping", and "photoshop contest") although Adobe disapproves of such use. Photoshop can edit and compose raster images in multiple layers and supports masks, alpha compositing and several color models. Photoshop uses its own PSD and PSB file formats to support these features. In addition to raster graphics, Photoshop has limited abilities to edit or render text and vector graphics (especially through clipping path for the latter), as well as 3D graphics and video. Its feature set can be expanded by plug-ins; programs developed and distributed independently of Photoshop that run inside it and offer new or enhanced features. Photoshop's naming scheme was initially based on version numbers. However, in October 2002 (following the introduction of Creative Suite branding), each new version of Photoshop was designated with "CS" plus a number; e.g., the eighth major version of Photoshop was Photoshop CS and the ninth was Photoshop CS2. Photoshop CS3 through CS6 were also distributed in two different editions: Standard and Extended. With the introduction of the Creative Cloud branding in June 2013 (and in turn, the change of the "CS" suffix to "CC"), Photoshop's licensing scheme was changed to that of software as a service subscription model. Historically, Photoshop was bundled with additional software such as Adobe ImageReady, Adobe Fireworks, Adobe Bridge, Adobe Device Central and Adobe Camera RAW. Alongside Photoshop, Adobe also develops and publishes Photoshop Elements, Photoshop Lightroom, Photoshop Express, Photoshop Fix, Adobe Illustrator, and Photoshop Mix. As of November 2019, Adobe has also released a full version of Photoshop for the iPad, and while initially limited, Adobe plans to bring more features to Photoshop for iPad. Collectively, they are branded as "The Adobe Photoshop Family". Early history Photoshop was developed in 1987 by two brothers, Thomas and John Knoll, who sold the distribution license to Adobe Systems Incorporated in 1988. Thomas Knoll, a Ph.D. student at the University of Michigan, began writing a program on his Macintosh Plus to display grayscale images on a monochrome display. This program (at that time called Display) caught the attention of his brother John, an Industrial Light & Magic employee, who recommended that Thomas turn it into a full-fledged image editing program. Thomas took a six-month break from his studies in 1988 to collaborate with his brother on the program. Thomas renamed the program ImagePro, but the name was already taken. Later that year, Thomas renamed his program Photoshop and worked out a short-term deal with scanner manufacturer Barneyscan to distribute copies of the program with a slide scanner; a "total of about 200 copies of Photoshop were shipped" this way. During this time, John traveled to Silicon Valley and gave a demonstration of the program to engineers at Apple Computer and Russell Brown, art director at Adobe. Both showings were successful, and Adobe decided to purchase the license to distribute in September 1988. While John worked on plug-ins in California, Thomas remained in Ann Arbor writing code. Photoshop 1.0 was released on February 19, 1990, for Macintosh exclusively. The Barneyscan version included advanced color editing features that were stripped from the first Adobe shipped version. The handling of color slowly improved with each release from Adobe and Photoshop quickly became the industry standard in digital color editing. When Photoshop 1.0 was released, digital retouching on dedicated high-end systems (such as the Scitex) cost around $300 an hour for basic photo retouching. The list price of Photoshop 1.0 for Macintosh in 1990 was $895. Photoshop was initially only available on Macintosh. In 1993, Adobe chief architect Seetharaman Narayanan ported Photoshop to Microsoft Windows. The Windows port led to Photoshop reaching a wider mass market audience as Microsoft's global reach expanded within the next few years. On March 31, 1995, Adobe purchased the rights for Photoshop from Thomas and John Knoll for $34.5 million so Adobe would no longer need to pay a royalty for each copy sold. File format Photoshop files have default file extension as .PSD, which stands for "Photoshop Document". A PSD file stores an image with support for all features of Photoshop; these include layers with masks, transparency, text, alpha channels and spot colors, clipping paths, and duotone settings. This is in contrast to many other file formats (e.g., .JPG or .GIF) that restrict content to provide streamlined, predictable functionality. A PSD file has a maximum height and width of 30,000 pixels, and a size limit of two gigabytes. From the beginning, Photoshop could save files in other formats, including TIF, JPEG, and GIF. These files are smaller than PSD files because they lack the editable features of a PSD file. These formats are required to use the file in publications or on the web. Adobe's discontinued program PageMaker required TIF format. Photoshop can also create and use files with the extension .PSB, which stands for "Photoshop Big" (also known as "large document format"). A PSB file extends the PSD file format, increasing the maximum height and width to 300,000 pixels and the size limit to around 4 exabytes. PSD and PSB formats are documented. Because of Photoshop's popularity, PSD files are widely used and supported to some extent by most competing software, including GIMP, Affinity Photo, and Clip Studio Paint. The .PSD file format can be exported to and from Adobe's other apps, such as Adobe Illustrator, Adobe Premiere Pro, and After Effects. Plugins Photoshop functionality can be extended by add-on programs called Photoshop plugins (or plug-ins). Adobe creates some, such as Adobe Camera Raw, but most are developed by third-parties. Some are free and some are commercial software. Most plugins work with only Photoshop or Photoshop-compatible hosts, but a few can also be run as standalone applications. There are various types of plugins, such as filter, export, import, selection, color correction, and automation. The most popular plugins are the filter plugins (also known as a 8bf plugins), available under the Filter menu in Photoshop. Filter plugins can either modify the current image or create content. Below are some popular types of plugins, and some well-known companies associated with them: Color correction plugins (Alien Skin Software, Nik Software, OnOne Software, Topaz Labs Software, The Plugin Site, etc.) Special effects plugins (Alien Skin Software, Auto FX Software, AV Bros., Flaming Pear Software, etc.) 3D effects plugins (Andromeda Software, Strata, etc.) Adobe Camera Raw (also known as ACR and Camera Raw) is a special plugin, supplied free by Adobe, used primarily to read and process raw image files so that the resulting images can be processed by Photoshop. It can also be used from within Adobe Bridge. Cultural impact Photoshop and derivatives such as Photoshopped (or just Shopped) have become verbs that are sometimes used to refer to images edited by Photoshop, or any image manipulation program. The same happens not only in English but as the Portuguese Wikipedia entry for image manipulation attests, even in that language, with the trademark being followed by the Portuguese verb termination -ar, yielding the word "photoshopar" (to photoshop). Such derivatives are discouraged by Adobe because, in order to maintain validity and protect the trademark from becoming generic, trademarks must be used as proper nouns. Version history Photoshop's naming scheme was initially based on version numbers, from version 0.63 (codename "Bond"; double-oh-seven), through version 0.87 (codename "Seurat" which was the first commercial version, sold as "Barneyscan XP"), version 1.0 (February 1990) all the way to version 7.0.1. Adobe published 7 major and many minor versions before the October 2003 introduction of version 8.0 which brought with it the Creative Suite branding. In February 2013 Adobe donated the source code of the 1990 1.0.1 version of Photoshop to the Computer History Museum. Pre-release versions Version 0.63 (October 1988) was the first known copy of Photoshop, though it was never publicly released. Version 0.87 (March 1989) was the first publicly available version of Photoshop, distributed commercially under the name "Barneyscan XP". Notable milestone features would be: Filters, Colour Separation, Virtual Memory (1.0), Paths, CMYK color (2.0), 16-bits-per-channel support, availability on Microsoft Windows (2.5), Layers, tabbed Palettes (3.0), Adjustments, Actions, Freeform Transform, PNG support (4.0), Editable Type, Magnetic Lasso and Pen, Freeform Pen, Multiple Undo, Layer Effects (5.0), Save For Web (5.5), Vector Shapes, revised User Interface (6.0), Vector Text, Healing Brush, Spell Check (7.0), Camera RAW (7.0.1). Version 1 Photoshop 1.0 was released in February 1990. Version 2 Photoshop 2.0 was released in June 1991. It added support for paths and the CMYK color model. Photoshop 2.5, released in November 1992, was the first version available for Windows. Version 3 Photoshop 3.0 was released in September 1994 for Mac OS 7. The Windows version came out later in November. Notably, this was the first version that brought Layers. Version 4 Photoshop 4.0 was released in November 1996. Version 5 Photoshop 5.0 was released in May 1998. Version 6 Photoshop 6.0 was released in September 2000. Version 7 Photoshop 7.0 was released in March 2002. CS (version 8) The first Photoshop CS was commercially released in October 2003 as the eighth major version of Photoshop. Photoshop CS increased user control with a reworked file browser augmenting search versatility, sorting and sharing capabilities and the Histogram Palette which monitors changes in the image as they are made to the document. Match Color was also introduced in CS, which reads color data to achieve a uniform expression throughout a series of pictures. CS2 (version 9) Photoshop CS2, released in May 2005, expanded on its predecessor with a new set of tools and features. It included an upgraded Spot Healing Brush, which is mainly used for handling common photographic problems such as blemishes, red-eye, noise, blurring and lens distortion. One of the most significant inclusions in CS2 was the implementation of Smart Objects, which allows users to scale and transform images and vector illustrations without losing image quality, as well as create linked duplicates of embedded graphics so that a single edit updates across multiple iterations. Adobe responded to feedback from the professional media industry by implementing non-destructive editing as well as the producing and modifying of 32-Bit High Dynamic Range (HDR) images, which are optimal for 3D rendering and advanced compositing. FireWire Previews could also be viewed on a monitor via a direct export feature. Photoshop CS2 brought the Vanishing Point and Image Warping tools. Vanishing Point makes tedious graphic and photo retouching endeavors much simpler by letting users clone, paint and transform image objects while maintaining visual perspective. Image Warping makes it easy to digitally distort an image into a shape by choosing on-demand presets or by dragging control points. The File Browser was upgraded to Adobe Bridge, which functioned as a hub for productivity, imagery and creativity, providing multi-view file browsing and smooth cross-product integration across Adobe Creative Suite 2 software. Adobe Bridge also provided access to Adobe Stock Photos, a new stock photography service that offered users one-stop shopping across five elite stock image providers to deliver high-quality, royalty-free images for layout and design. Camera Raw version 3.0 was a new addition in CS2, and it allowed settings for multiple raw files to be modified simultaneously. In addition, processing multiple raw files to other formats including JPEG, TIFF, DNG or PSD, could be done in the background without executing Photoshop itself. Photoshop CS2 brought a streamlined interface, making it easier to access features for specific instances. In CS2 users were also given the ability to create their own custom presets, which was meant to save time and increase productivity. In January 2013, Photoshop CS2 was released with a published serial number due to a technical glitch in Adobe's CS2 activation servers (see Creative Suite 1 and 2). CS3 (version 10) CS3 and CS3 Extended were released in April 2007 to the United States and Canada. They were also made available through Adobe's online store and Adobe Authorized Resellers. Both CS3 and CS3 Extended are offered as either a stand-alone application or feature of Adobe Creative Suite. Both products are compatible with Intel-based Macs and PowerPCs, supporting Windows XP and Windows Vista. CS3 is the first release of Photoshop that will run natively on Macs with Intel processors: previous versions can only run through the translation layer Rosetta, and will not run at all on Macs running Mac OS X 10.7 or later. CS3 improves on features from previous versions of Photoshop and introduces new tools. One of the most significant is the streamlined interface which allows increased performance, speed, and efficiency. There is also improved support for Camera RAW files which allow users to process images with higher speed and conversion quality. CS3 supports over 150 RAW formats as well as JPEG, TIFF and PDF. Enhancements were made to the Black and White Conversion, Brightness and Contrast Adjustment and Vanishing Point Module tools. The Black and White adjustment option improves control over manual grayscale conversions with a dialog box similar to that of Channel Mixer. There is more control over print options and better management with Adobe Bridge. The Clone Source palette is introduced, adding more options to the clone stamp tool. Other features include the nondestructive Smart Filters, optimizing graphics for mobile devices, Fill Light and Dust Busting tools. Compositing is assisted with Photoshop's new Quick Selection and Refine Edge tools and improved image stitching technology. CS3 Extended includes everything in CS3 and additional features. There are tools for 3D graphic file formats, video enhancement and animation, and comprehensive image measurement and analysis tools with DICOM file support. The 3D graphic formats allow 3D content to be incorporated into 2D compositions. As for video editing, CS3 supports layers and video formatting so users can edit video files per frame. CS4 (version 11) CS4 and CS4 Extended were released on October 15, 2008. They were also made available through Adobe's online store and Adobe Authorized Resellers. Both CS4 and CS4 Extended are offered as either a stand-alone application or feature of Adobe Creative Suite. Both products are compatible with Intel-based Mac OS X and PowerPCs, supporting Windows XP and Windows Vista. CS4 features smoother panning and zooming, allowing faster image editing at a high magnification. The interface is more simplified with its tab-based interface making it cleaner to work with. Photoshop CS4 features a new 3D engine allowing the conversion of gradient maps to 3D objects, adding depth to layers and text, and getting print-quality output with the new ray-tracing rendering engine. It supports common 3D formats; the new Adjustment and Mask panels; content-aware scaling (seam carving); fluid canvas rotation and File display options. The content-aware scaling allows users to intelligently size and scale images, and the canvas rotation tool makes it easier to rotate and edit images from any angle. Adobe released Photoshop CS4 Extended, which has the features of Adobe Photoshop CS4, plus capabilities for scientific imaging, 3D, motion graphics, accurate image analysis and high-end film and video users. The faster 3D engine allows users to paint directly on 3D models, wrap 2D images around 3D shapes and animate 3D objects. As the successor to Photoshop CS3, Photoshop CS4 is the first x64 edition of Photoshop on consumer computers for Windows. The color correction tool has also been improved significantly. CS5 (version 12) Photoshop CS5 was launched on April 12, 2010. In a video posted on its official Facebook page, the development team revealed the new technologies under development, including three-dimensional brushes and warping tools. In May 2011, Adobe Creative Suite 5.5 (CS5.5) was released, with new versions of some of the applications. Its version of Photoshop, 12.1, is identical to the concurrently released update for Photoshop CS5, version 12.0.4, except for support for the new subscription pricing that was introduced with CS5.5. CS5 introduces new tools such as the Content-Aware Fill, Refine Edge, Mixer Brush, Bristle Tips and Puppet Warp. The community also had a hand in the additions made to CS5 as 30 new features and improvements were included by request. These include automatic image straightening, the Rule-of-Thirds cropping tool, color pickup, and saving a 16-bit image as a JPEG. Another feature includes the Adobe Mini Bridge, which allows for efficient file browsing and management. CS5 Extended includes everything in CS5 plus features in 3D and video editing. A new materials library was added, providing more options such as Chrome, Glass, and Cork. The new Shadow Catcher tool can be used to further enhance 3D objects. For motion graphics, the tools can be applied to over more than one frame in a video sequence. CS5 and CS5 Extended were made available through Adobe's online store, Adobe Authorized Resellers and Adobe direct sales. Both CS5 and CS5 Extended are offered as either a stand-alone application or a feature of Adobe Creative Suite 5. Both products are compatible with Intel-based Mac OS X and Windows XP, Windows Vista, and Windows 7. CS6 (version 13) Photoshop CS6, released in May 2012, added new creative design tools and provided a redesigned interface with a focus on enhanced performance. New features have been added to the Content-Aware tool such as the Content-Aware Patch and Content-Aware Move. Adobe Photoshop CS6 brought a suite of tools for video editing. Color and exposure adjustments, as well as layers, are among a few things that are featured in this new editor. Upon completion of editing, the user is presented with a handful of options of exporting into a few popular formats. CS6 brings the "straighten" tool to Photoshop, where a user simply draws a line anywhere on an image, and the canvas will reorient itself so that the line drawn becomes horizontal, and adjusts the media accordingly. This was created with the intention that users will draw a line parallel to a plane in the image, and reorient the image to that plane to more easily achieve certain perspectives. CS6 allows background saving, which means that while another document is compiling and archiving itself, it is possible to simultaneously edit an image. CS6 also features a customizable auto-save feature, preventing any work from being lost. With version 13.1.3, Adobe dropped support for Windows XP (including Windows XP Professional x64 Edition); thus, the last version that works on Windows XP is 13.0.1. Adobe also announced that CS6 will be the last suite sold with perpetual licenses in favor of the new Creative Cloud subscriptions, though they will continue to provide OS compatibility support as well as bug fixes and security updates as necessary. Starting January 9, 2017, CS6 is no longer available for purchase, making a Creative Cloud license the only purchase option going forward. No more updates will be available for all CS6 software either. CC (version 14) Photoshop CC (14.0) was launched on June 18, 2013. As the next major version after CS6, it is only available as part of a Creative Cloud subscription. Major features in this version include new Smart Sharpen, Intelligent Upsampling, and Camera Shake Reduction for reducing blur caused by camera shake. Editable Rounded Rectangles and an update to Adobe Camera Raw (8.0) were also included. Since the initial launch, Adobe has released two additional feature-bearing updates. The first, version 14.1, was launched on September 9, 2013. The major features in this version were Adobe Generator, a Node.js-based platform for creating plug-ins for Photoshop. Photoshop 14.1 shipped with two plug-ins, one to automatically generate image assets based on an extension in the layer name, and another to automatically generate assets for Adobe Edge Reflow. Version 14.2 was released on January 15, 2014. Major features include Perspective Warp, Linked Smart Objects, and 3D Printing support. CC 2014 (version 15) Photoshop CC 2014 (15.0) was released on June 18, 2014. CC 2014 features improvements to content-aware tools, two new blur tools (spin blur and path blur) and a new focus mask feature that enables the user to select parts of an image based on whether they are in focus or not. Other minor improvements have been made, including speed increases for certain tasks. CC 2015 (version 16 and version 17) Photoshop CC 2015 was released on June 15, 2015. Adobe added various creative features including Adobe Stock, which is a library of custom stock images. It also includes and have the ability to have more than one layer style. For example, in the older versions of Photoshop, only one shadow could be used for a layer but in CC 2015, up to ten are available. Other minor features like Export As, which is a form of the Save For Web in CC 2014 were also added. The updated UI as of November 30, 2015, delivers a cleaner and more consistent look throughout Photoshop, and the user can quickly perform common tasks using a new set of gestures on touch-enabled devices like Microsoft Surface Pro. CC 2015 also marks the 25th anniversary of Photoshop. CC 2017 (version 18) Photoshop CC 2017 was released on November 2, 2016. It introduced a new template selector when creating new documents, the ability to search for tools, panels and help articles for Photoshop, support for SVG OpenType fonts and other small improvements. In December 2016, a minor update was released to include support for the MacBook Pro Touch Bar. CC 2018 (version 19) Photoshop CC 2018 (version 19) was released on October 18, 2017. It featured an overhaul to the brush organization system, allowing for more properties (such as color and opacity) to be saved per-brush and for brushes to be categorized in folders and sub-folders. It also added brush stroke smoothing, and over 1000 brushes created by Kyle T. Webster (following Adobe's acquisition of his website, KyleBrush.com). A Curvature Pen tool, similar to the one in Illustrator, was added, allowing for faster creation of Bézier paths. Other additions were Lightroom Photo access, Variable font support, select subject, copy-paste layers, enhanced tooltips, 360 panorama and HEIF support, PNG compression, increased maximum zoom level, symmetry mode, algorithm improvements to Face-aware and selection tools, color and luminance range masking, improved image resizing, and performance improvements to file opening, filters, and brush strokes. CC 2019 (version 20) Photoshop CC 2019 was released on October 15, 2018. Beginning with Photoshop CC 2019 (version 20.0), the 32-bit version of Windows is no longer supported. This version Introduced a new tool called Frame Tool to create placeholder frames for images. It also added multiple undo mode, auto-commitment, and prevented accidental panel moves with lock work-space. Live blend mode previews are added, allowing for faster scrolling over different blend mode options in the layers panel. Other additions were Color Wheel, Transform proportionally without Shift key, Distribute spacing like in Illustrator, ability to see longer layer names, match font with Japanese fonts, flip document view, scale UI to font, reference point hidden by default, new compositing engine, which provides a more modern compositing architecture is added which is easier to optimize on all platforms. 2020 (version 21) Photoshop 2020 was released on November 4, 2019. Version 21 has many new and enhanced features like the new object selection tool for better automate complex selections, new properties panel, enhanced transform warp, new keyboard shortcuts for paint & brush and background image removal option. It added several improvements to the new content-aware fill and to the new document tab. Also added were animated GIF support, improved lens blur performance and one-click zoom to a layer's contents. It introduced new swatches, gradients, patterns, shapes and stylistic sets for OpenType fonts. With this version users now can easily convert smart objects to layers and also can adjust 32-bit layers for brightness/contrast and curves. Presets are now more intuitive to use and easier to organize. With the February 2020 update (version 21.1) Photoshop now can iteratively fill multiple areas of an image without having to leave content-aware fill workspace. This version improved GPU based lens blur quality and provided performance improvements, such as accelerating workflows with smoother panning, zooming and navigation of documents. Version 21 was the first version where the iPad version was released. With Photoshop on the iPad, combined with the new Cloud PSD file format, a user can save cloud documents and work across Windows, Mac and iPad. Photoshop on the iPad does not have all the features of the desktop Photoshop. Adobe promises to update Photoshop on the iPad at "a much more aggressive pace than it has with its current Creative Cloud apps for the desktop". Adobe has provided a timeline for enhancing Photoshop on the iPad to have more of the features of desktop Photoshop. Version 21.2 of the desktop version was released in June 2020. It introduced faster portrait selection, Adobe Camera Raw improvements, auto-activated Adobe Fonts, rotatable patterns, and improved Match Font. 2021 (version 22) Version 22.0.0 was released in October 2020. Version 22.0.1 was released in November 2020. Version 22.1.0 was released in December 2020. Version 22.1.1 was released in January 2021. Version 22.2 was released in February 2021. Version 22.3 was released in March 2021. This is the first macOS release to run natively on Apple silicon. Version 22.3.1 was released in April 2021. Version 22.4 was released in May 2021. Version 22.4.1 was released in May 2021. Version 22.4.2 was released in June 2021. Version 22.4.3 was released in July 2021. Version 22.5 was released in August 2021. Version 22.5.1 was released in September 2021. Final version to include (as Legacy Swatches) a large set of Pantone color books, including Pastels & Neons, and Premium Metallics. 2022 (version 23) Version 23.0 was released in October 2021. First version to include (as Legacy Swatches) a reduced set of Pantone color books - only CMYK, Metallics and Solid colors. Content Credentials (Beta) was introduced. When enabled, the editing information is captured in a tamper-evident form and resides with the file through successive copy generations. It aligns with the C2PA standard on digital provenance across the internet. Version 23.0.1 was released in November 2021. Version 23.0.2 was released in November 2021. Version 23.1 was released in December 2021. Version 23.1.1 was released in January 2022. Version 23.2 was released in February 2022. Version 23.3 was released in April 2022. Version 23.4 was released in June 2022. Version 23.5 was released in August 2022. Final version with built-in support for basic Pantone colors - CMYK, Metallics and Solid (as Legacy Swatches). Future versions require a separate subscription to access Pantone colors. 2023 (version 24) Version 24.0 was released in October 2022. First version without built-in support for Pantone colors. All Pantone colors have been removed as of August 16, 2022. Separate paid subscription now required to access Pantone Connect extension and download color "fandecks". Version 24.1 was released in December 2022. Version 24.4.1 was released on April 20, 2023. Version 24.7 was released on July 27, 2023. 2024 (version 25) Version 25.0 was released in September 2023. This version added Generative Fill and Generative Expand for commercial use. Adobe Photoshop family The Adobe Photoshop family is a group of applications and services made by Adobe for the use of professional image editing. Several features of the Adobe Photoshop family are pixel manipulating, image organizing, photo retouching, and more. Current applications Bridge is an image organizer and digital asset management app. It features limited integration with other Adobe apps but has no editing capabilities of its own. DNG Converter is a tool used to convert DNG files into other file formats. Elements Organizer is the digital asset management app for Photoshop Elements and Premiere Elements. It is able to organize photos and video projects in one place. Fresco is a mobile drawing and painting app, developed initially for iOS and marketed by Adobe through Creative Cloud. It was later adapted to run on certain Windows devices and Microsoft Surface tablets. Photoshop Lightroom is a creative image organization and image manipulation software developed by Adobe Inc. as part of the Creative Cloud subscription family. Lightroom Classic is the offline desktop version of the Photoshop Lightroom photo editing and viewing applications offered by subscription through Creative Cloud. Photoshop Camera is an image tool that easily captures and shares photos with your camera. Photoshop Elements is a graphics editor for photographers, image editors and hobbyists. It contains most of the features of the professional version but with fewer and simpler options. The program allows users to create, edit, organize and share images. Originally introduced alongside Adobe Photoshop version 6, Photoshop Elements targets photography enthusiasts and thus lacks many features that make it useful in a proper print production environment. Photoshop Elements is available for Windows and macOS. It is not available as part of a Creative Cloud subscription, but rather as a single purchase or upgrade purchase. Photoshop Express is an image editing and collage making mobile application from Adobe Inc. The app is available on iOS, Android and Windows phones and tablets. It can also be installed on Windows desktop via the Microsoft Store. Photoshop Express Editor has various features which can be used to enhance photos. In November 2016 Collage creation was introduced to Adobe Photoshop Express on iOS. They allow editing pictures in the smartphone or tablet rather than online. It can be used to showcase your latest art, ideas, or products. In the early days of the beta test, the product's terms of use raised controversy, in that it claimed to retain an irrevocable license to use certain works submitted by end users in perpetuity. At that time, Adobe held users of Photoshop Express to Adobe.com's general Terms of Use, which had not been drafted in contemplation of the sort of user-created content utilized in Photoshop Express. Following user concerns and negative press, Adobe issued new, more specialized Terms of Use for the Photoshop Express product that superseded sections of the General Terms, and clarified many of these issues. Changes included making the license expressly revocable and indicating that Adobe's rights to use the content are solely for the operation of Photoshop Express itself. Photoshop Fix is a photo retouching app for mobile devices marketed through Creative Cloud. The retouched images can be exported to the desktop version of Photoshop for further work. Photoshop Mix is a mobile application designed as a replacement for Photoshop Touch specifically for tablets and touchscreen devices. It includes many of the features of the personal computer version, including layers, selection tools, adjustments, and filters. Edited files could be synced with Creative Cloud. Photoshop Sketch is a drawing and painting app for mobile devices marketed through Creative Cloud. Sketches made can be exported to Photoshop. They can also be uploaded directly to the Bēhance social media platform. Discontinued applications ActiveShare is a discontinued photo-sharing platform distributed by Adobe Systems. The Photoshop Album application replaced ActiveShare in 2003. Fireworks is a discontinued raster graphics editor for web designers. It could create interactive contents (e.g. buttons that change shape when the mouse cursor is hovered on) and animations. ImageReady is a discontinued raster graphics editor for web designers. It was discontinued on CS2 in favor of Fireworks. Photoshop Album is a piece of application software designed to import, organize and edit digital photos, and allows quick and easy searching and sharing of entire photo collections. It was initially released on February 18, 2003. The last version was Photoshop Album 3.2.0. It was discontinued in favor of Photoshop Elements. It has been compared to Apple Inc.'s iPhoto and Google's Picasa. Photoshop Limited Edition (LE) was a graphics editor for novice photographers and hobbyists. It contains most of the features of the professional version but with fewer and simpler options. It was instead replaced by Photoshop Elements in September 2000. Preview CC is an app for previewing mobile designs.	Technology	Multimedia_2	null
24900	https://en.wikipedia.org/wiki/Plastic%20explosive	Plastic explosive	Plastic explosive is a soft and hand-moldable solid form of explosive material. Within the field of explosives engineering, plastic explosives are also known as putty explosives or blastics. Plastic explosives are especially suited for explosive demolition. Common plastic explosives include Semtex and C-4. The first manufactured plastic explosive was gelignite in 1875, invented by Alfred Nobel. Usage Plastic explosives are especially suited for explosive demolition of obstacles and fortifications by combat engineers as they can be easily formed into ideal shapes for cutting structural members and have a high enough velocity of detonation and density for metal cutting work. An early use of plastic explosives was in the warhead of the Petard demolition mortar of the British Armoured Vehicle Royal Engineers (AVRE) which was used to destroy concrete fortifications encountered during Operation Overlord (D-Day). The original use of Nobel 808 supplied by the SOE was for sabotage of German installations and railways in Occupied Europe. They are generally not used for ordinary blasting as they tend to be significantly more expensive than other materials that perform just as well. A common commercial use of plastic explosives is for shock hardening high manganese percentage steel, a material typically used for train rail components and earth digging implements. Reactive armor in tanks uses plastic explosives sandwiched between two plates of steel. Incoming high explosive shaped charge anti-tank rounds pierce the outer steel plate, then detonate the plastic explosive. This disrupts the energy from the incoming round and shields the tank. History The first plastic explosive was gelignite, invented by Alfred Nobel in 1875. Prior to World War I, the British explosives chemist Oswald Silberrad obtained British and U.S. patents for a series of plastic explosives called "Nitrols", composed of nitrated aromatics, collodion, and oxidising inorganic salts. The language of the patents indicate that at this time, Silberrad saw no need to explain to "those versed in the art" either what he meant by plasticity or why it may be advantageous, as he only explains why his plastic explosive is superior to others of that type. One of the simplest plastic explosives was Nobel's Explosive No. 808, of the gelignite type, also known as Nobel 808 (often just called Explosive 808 in the British Armed Forces during the Second World War), developed by the British company Nobel Chemicals Ltd well before World War II. It had the appearance of green plasticine with a distinctive smell of almonds. During World War II it was extensively used by the British Special Operations Executive (SOE) at Aston House for sabotage missions. It is also the explosive used in HESH anti-tank shells and was an essential factor in the devising of the Gammon grenade. Captured SOE-supplied Nobel 808 was the explosive used in the failed 20 July plot assassination attempt on Adolf Hitler in 1944. During and after World War II a number of new RDX-based explosives were developed, including Compositions C, C2, and eventually C3. Together with RDX, these incorporate various plasticizers to decrease sensitivity and make the composition plastic. The origin of the obsolete term "plastique" dates back to the Nobel 808 explosive introduced to the U.S. by the British in 1940. The samples of explosive brought to the U.S. by the Tizard Mission had already been packaged by the SOE ready for dropping via parachute container to the French Resistance and were therefore labeled in French, as Explosif Plastique. It is still referred to by this name in France and also by some Americans. Types Composition C The British used a plastic explosive during World War II as a demolition charge. The specific explosive, Composition C, was 88.3% RDX and 11.7% non-oily, non-explosive plasticizer. The material was plastic between , but was brittle at colder temperatures and gummy at higher temperatures. Composition C was superseded by Composition C2, which used a mixture of 80% RDX and 20% plasticizer. Composition C2 had a wider temperature range at which it remained plastic, from . Composition C2 was replaced by Composition C3, which was a mixture of 77% RDX and 23% explosive plasticizer. C3 was effective but proved to be too brittle in cold weather and was replaced with C4. There are three classes of C4, with varying amounts of RDX and polyisobutylene. Semtex List of plastic explosives Australia: PE4, PE4-MC Austria: KNAUERIT SPEZIAL Czech Republic: Semtex-1H (orange-colored), Semtex 1A (red-colored), Semtex 10 (also called Pl Np 10; black-colored), Pl Hx 30 (gray-colored) Finland: PENO France: Hexomax, Composition C-4 PLASTRITE (FORMEX P1, Pla Np 87) Germany: Sprengkörper DM12, P8301, Seismoplast 1 (Sprengmasse, formbar) Netherlands: Knaverit S1 (light orange-colored) Greece: C3, C4 India: PEK-1 Israel: Semtex Italy: T-4 Plastico Norway: NM91 (HMX), C4, DPX10 (PE8) Pakistan: PE-3A Poland: PMW, NITROLIT Russia: PVV-5A Plastic Explosive Slovakia: CHEMEX (Composition C-4 equivalent), TVAREX 4A, Pl Hx 30 South Africa: PE9 (Composition C-4 equivalent) Spain: PG2, PG4, GOMA 0, GOMA 1, GOMA 2 Sweden: Sprängdeg m/46, NSP711 (PETN-based), NSH711 (cyclonite-based) Switzerland: PLASTEX produced by SSE Turkey: Composition C-4 United Kingdom MOD (Ministry of Defence) explosives: PE2 (sheet explosive, superseded by SX2), PE3A (superseded by PE4), PE4 (pure to off-white slab, block, or stick, superseded by PE7 and PE8 in MOD usage), SX2 (sheet explosive, superseded by SX4), PE7 (pure to off-white slab or block, Hexomax variant), PE8 (pure to off-white slab or block, current in-service slab charge), SX4 (sheet explosive), DPX (DPX1 used in L26A1 Bangalore Torpedo Demolition Charge, DPX9 used in SABREX and as a key component of SX4) Non-MOD explosives: Composition C-4 (M5A1 and M112 charges produced by Mondial Defence Systems), Semtex (Several variants including Razor produced by Mondial Defence Systems, PW4 variant produced by Chemring)) United States: Composition C-4 (pure white block or sheet, current in-service charges designated as M112 and M118), PETN and RDX based Sheet Explosive (Primasheet, Durasheet), DURABLOCK Advanced Demolition Explosive(ADX) Produced by Donovan Commercial Industries Yugoslavia/Serbia: PP–01 (Composition C-4 equivalent)	Technology	Explosive weapons	null
11406745	https://en.wikipedia.org/wiki/Delichon	Delichon	Delichon is a small genus of passerine birds that belongs to the swallow family and contains four species called house martins. These are chunky, bull-headed and short-tailed birds, blackish-blue above with a contrasting white rump, and with white or grey underparts. They have feathering on the toes and tarsi that is characteristic of this genus. The house martins are closely related to other swallows that build mud nests, particularly the Hirundo barn swallows. They breed only in Europe, Asia and the mountains of North Africa. Three species, the common, Siberian and Asian house martins, migrate south in winter, while the Nepal house martin is resident in the Himalayas year-round. The house martins nest in colonies on cliffs or buildings, constructing feather- or grass-lined mud nests. The typical clutch is two or three white eggs; both parents build the nest, incubate the eggs and feed the chicks. These martins are aerial hunters of small insects such as flies and aphids. Despite their flying skills the Delichon martins are sometimes caught by fast-flying birds of prey. They may carry fleas or internal parasites. None of the species are considered threatened, although widespread reductions in common house martin numbers have been reported from central and northern Europe. This decline is due to factors including poor weather, poisoning by agricultural pesticides, lack of mud for nest building and competition with house sparrows for nest sites. Taxonomy The four Delichon species are members of the swallow family of birds, and are classed as members of the Hirundininae subfamily which comprises all swallows and martins except the very distinctive river martins. DNA studies suggest that there are three major groupings within the Hirundininae, broadly correlating with the type of nest built. The groups are the "core martins" including burrowing species like the sand martin, the "nest-adopters", which are birds like the tree swallow that utilise natural cavities, and the "mud nest builders". The Delichon species construct a closed mud nest and therefore belong to the latter group; they appear to be intermediate between the Hirundo and Ptyonoprogne species that make open cup nests, and the Cecropis and Petrochelidon swallows, which have retort-like closed nests with an entrance tunnel. The genetic evidence suggests a close relationship between Hirundo and Delichon, which is further supported by the frequency of interbreeding between two widespread species, the barn swallow and the common house martin, despite being their being in different genera. The suggested taxonomic sequence of the mud-building swallows has been recommended by at least two European taxonomic committees. The genus Delichon was created by American naturalist Thomas Horsfield and British entomologist Frederic Moore in 1854 to accommodate the Nepal house martin that was first described by Moore in the same year, and is therefore the type species for the genus. The two other house martins were moved to Delichon from the genus Chelidon in which they had been placed up to that time. In 2021, the Siberian house martin was split from the common house martin based on morphological and vocal differences. The name of the genus, "Delichon", is an anagram of the Ancient Greek term χελιδὡν/chelidôn, meaning swallow. Species The genus contains four similar species: The common and Asian house martins have sometimes been considered to be a single species, although both breed in the western Himalayas without hybridising. There is also limited DNA evidence that suggests a significant genetic distance between these two martins. Distribution and habitat Delichon is an Old World genus with all four species breeding only in the Northern Hemisphere. The common house martin is a widespread migrant breeder across Europe, north Africa and all northern temperate Asia to Kamchatka. It winters in tropical Africa. The Siberian house martin breeds in northeast Russia and winters in southern Asia. The Asian house martin breeds further south than the Siberian house martin in the mountains of central and eastern Asia; its nominate subspecies winters in Southeast Asia, but the races breeding in the Himalayas and Taiwan may just move from the high mountains to lower altitudes. The Nepal house martin is resident in the mountains of southern Asia. The preferred habitat of the common and Siberian house martins is open country with low vegetation, such as pasture, meadows and farmland, and preferably near water, although it is also found in mountains up to at least 2,200 metres (7,200 ft) altitude. As the name implies, they readily nest on man-made buildings, and will breed even in city centres if the air is clean enough. The other two species favour mountainous country (and sea cliffs in the case of Asian house martin); they use buildings as nest sites less frequently than their northern relative. The wintering grounds of the two migrant species include a range of open country and hilly habitats. Description Delichon martins are 13–15 cm (5–6 in) long, blackish blue above with a contrasting white rump, and with white or grey underparts. They are chunky, bull-headed and short-tailed birds, and have feathering on the toes and tarsi. The common house martin is the largest bird, with an average weight of 18.3 g (0.65 oz), and has the most deeply forked tail; the Nepalese species is the smallest (15 g, 0.53 oz) and has the squarest tail. Distinctive species plumage features are the black chin and black undertail coverts of the Nepal house martin, and the greyish wash to the underparts of the Asian house martin. As with other swallows and martins, the moult is slow and protracted because of the need to maintain efficient flight at all times to enable feeding. Moult normally starts on arrival at the wintering grounds, but overlaps with the breeding season for the non-migratory Nepal house martin. The Delichon martins have simple flight calls of one to three notes. In the two more widespread species these have a distinctive buzzing quality. The male's song is a short simple ripple, perhaps less musical than that given by other swallows. As a group, the house martins cannot easily be confused with any other swallows. Four species of the genus Tachycineta have white rumps and underparts, but they have bright metallic green or blue-green upperparts, longer tails, and are restricted to Central and South America. The variable plumages of the South Asian species and a confused taxonomic history has left their distribution ranges in doubt. Behaviour Breeding The Delichon martins were originally cliff nesters, breeding in colonies situated under an overhang on a vertical cliff. However, the house martin now largely uses human structures, as, to a lesser extent, does the Asian house martin. The typical nest is a grass or feather-lined deep closed mud bowl with a small opening at the top, but many Asian house martins leave the top of the nest open. David Winkler and Frederick Sheldon believe that evolutionary development in the mud-building swallows, and individual species follow this order of construction. A retort builder like red-rumped swallow starts with an open cup, closes it, and then builds the entrance tunnel. Winkler and Sheldon propose that the development of closed nests reduced competition between males for copulations with the females. Since mating occurs inside the nest, the difficulty of access means other males are excluded. This reduction in competition permits the dense breeding colonies typical of the Delichon martins. The urban common house martin has to compete with house sparrows, which frequently attempt to take over the nest during construction, with the house martins rebuilding elsewhere if the sparrows are successful. The entrance at the top of the completed cup is so small that the sparrows cannot take over the nest once it is finished. As with other swallows, pairing and copulation displays are normally brief, taking just a few minutes. The male calls to a female and attempts to lead her to the nest, where he lands and continues calling while posing with lowered head, dropped wings and ruffled throat. If he is successful, the female calls and allows him to mount her, usually in the nest. Three or four white eggs are the normal clutch and all three species are frequently double-brooded. Both sexes build the nest, incubate the eggs and feed the chicks, although the female does most of the incubation, which normally lasts 14–16 days. The newly hatched chicks are altricial, and after a further 22–32 days, depending on weather, the chicks leave the nest. The fledged young stay with, and are fed by, the parents for about a week after leaving the nest. Occasionally, young birds from the first brood will assist in feeding the second brood. A Scottish study showed that mortality in common house martins occurred mostly outside the breeding season and averaged 57%. Females that had raised two clutches in a season had a higher mortality than those that were single-brooded, but there was no such correlation for the males. Feeding The Delichon species typically feed higher in the air, and take smaller prey than other swallows. It is believed that this reduces inter-specific food competition, particularly with the barn swallow which shares much of the breeding and wintering range of the martins. The insects eaten are mostly small flies, aphids and Hymenoptera such as winged ants. A wide range of other insects are caught, including Lepidoptera, beetles and lacewings. The Asian house martin appears to occasionally take terrestrial springtails and larvae and the common house martin also sometimes feeds on the ground. These martins are gregarious, feeding in flocks often with other aerial predators like swifts, or other hirundines such as the barn or striated swallows. In the case of at least the common house martin, the start of egg laying appears to be linked to the appearance of large numbers of flying aphids, which provide a stable and abundant food supply. Predators and parasites The main predators of the house martins are those birds of prey which are capable of catching these agile fliers, such as the hobby. Birds of the Delichon species are most vulnerable when collecting mud from the ground. This has therefore become a communal activity, with a group of these birds descending suddenly on a patch of mud. The usually insectivorous collared falconet has been recorded as hunting Nepal house martins. The house martins are parasitised by fleas and mites, including the "house martin flea", Ceratophyllus hirundinis and its relatives. A Polish study of the common house martin showed that nests typically contained more than 29 species of ectoparasite, with C. hirundinis and another swallow specialist, Oeciacus hirundinis, the most abundant. The genus also hosts endoparasites such as Haemoproteus prognei (avian malaria), which are transmitted by blood-sucking insects including mosquitoes. More than 40 beetle species have been recorded in common house martin nests, but most are either typical of the locality or found in the nests of other birds. The typical number of individuals, around 200, is relatively low compared to other bird species (1,400 individual beetles for house sparrow, 2,000 for sand martin). The beetles have no effect on the nesting birds, and the reason for their comparatively low numbers is unknown, although the numbers of specific parasites found in house martin's nests is also quite small. Conservation status The International Union for Conservation of Nature (IUCN) is the organisation responsible for assessing the conservation status of species. A species is assessed as subject to varying levels of threat if it has a small, fragmented or declining range, or if the total population is less than 10,000 mature individuals, or if numbers have dropped rapidly (by more than 10% in ten years or three generations). None of the Delichon species meets these criteria, and all four house martins are therefore considered of least concern. The numbers of the two southern Asian species are unknown, but both can be locally abundant, and the Asian house martin is extending its range in southern Siberia. The lowland breeding common house martin has greatly benefited from forest clearance, creating the open habitats it prefers, and from human habitation which has given it an abundance of safe man-made nest sites, although widespread declines in its numbers have been reported from central and northern Europe since 1970. This is due to factors including poor weather, poisoning by agricultural pesticides, lack of mud for nest building and competition with house sparrows for nest sites. The formerly conspecific Siberian house martin is also declining. Despite this, the huge geographical range and large numbers of the two northern house martins mean that their global status is secure. Fossil record Source: Delichon polgardiensis (late Miocene of Polgardi, Hungary) Delichon pusillus (Pliocene of Csarnota, Hungary) Delichon major (Pliocene of Beremend, Hungary) Citations	Biology and health sciences	Passerida	Animals
679350	https://en.wikipedia.org/wiki/Food%20allergy	Food allergy	A food allergy is an abnormal immune response to food. The symptoms of the allergic reaction may range from mild to severe. They may include itchiness, swelling of the tongue, vomiting, diarrhea, hives, trouble breathing, or low blood pressure. This typically occurs within minutes to several hours of exposure. When the symptoms are severe, it is known as anaphylaxis. A food intolerance and food poisoning are separate conditions, not due to an immune response. Common foods involved include cow's milk, peanuts, eggs, shellfish, fish, tree nuts, soy, wheat, and sesame. The common allergies vary depending on the country. Risk factors include a family history of allergies, vitamin D deficiency, obesity, and high levels of cleanliness. Allergies occur when immunoglobulin E (IgE), part of the body's immune system, binds to food molecules. A protein in the food is usually the problem. This triggers the release of inflammatory chemicals such as histamine. Diagnosis is usually based on a medical history, elimination diet, skin prick test, blood tests for food-specific IgE antibodies, or oral food challenge. Management involves avoiding the food in question and having a plan if exposure occurs. This plan may include giving adrenaline (epinephrine) and wearing medical alert jewelry. Early childhood exposure to potential allergens may be protective against later development of a food allergy. The benefits of allergen immunotherapy for treating food allergies are not proven, thus not recommended . Some types of food allergies among children resolve with age, including those to milk, eggs, and soy; while others such as to nuts and shellfish typically do not. In the developed world, about 4% to 8% of people have at least one food allergy. They are more common in children than adults and appear to be increasing in frequency. Male children appear to be more commonly affected than females. Some allergies more commonly develop early in life, while others typically develop in later life. In developed countries, more people believe they have food allergies when they actually do not have them. Signs and symptoms Food allergy symptoms occur within minutes to hours after exposure and may include: Rash Hives Itching of mouth, lips, tongue, throat, eyes, skin, or other areas Swelling (angioedema) of lips, tongue, eyelids, or the whole face Difficulty swallowing Runny or congested nose Hoarse voice Wheezing and/or shortness of breath Diarrhea, abdominal pain, and/or stomach cramps Lightheadedness Fainting Nausea Vomiting In some cases, however, onset of symptoms may be delayed for hours. Symptoms can vary. The amount of food needed to trigger a reaction also varies. Serious danger regarding allergies can begin when the respiratory tract or blood circulation is affected. The former can be indicated through wheezing and cyanosis. Poor blood circulation leads to a weak pulse, pale skin and fainting. A severe case of an allergic reaction, caused by symptoms affecting the respiratory tract and blood circulation, is called anaphylaxis. When symptoms are related to a drop in blood pressure, the person is said to be in anaphylactic shock. Anaphylaxis occurs when IgE antibodies are involved, and areas of the body that are not in direct contact with the food become affected and show symptoms. Those with asthma or an allergy to peanuts, tree nuts, or seafood are at greater risk for anaphylaxis. Causes Common food allergies Allergic reactions are abnormal immune responses that develop after exposure to a given food allergen. Food allergens account for about 90% of all allergic reactions. The most common food allergens include milk, eggs, peanuts, tree nuts, fish, shellfish, soy, and wheat, which are referred to as "the big eight", and are required by US law to be on labels of foods that contain those foods. In April 2021, President Biden signed the FASTER Act into law. This recognized sesame as the ninth US mandatory food label allergen. Peanuts, a member of the legume family, are one of the most common food allergens that induce reactions in both children and adults. Affecting about 2% of the Western population, peanut allergies tend to cause more severe reactions and anaphylaxis than other food allergies. Tree nuts, including almonds, brazil nuts, cashews, coconuts, hazelnuts, macadamia nuts, pecans, pistachios, pine nuts, and walnuts, are also common allergens. Affected individuals may be sensitive to one particular tree nut or many different ones. Peanuts and seeds, including sesame seeds and poppy seeds, can be processed to extract oils, but trace amounts of protein may also elicit an allergic reaction. Peanut and tree nut allergies are lifelong conditions for the majority of those affected, although evidence shows that ~20% of those with peanut allergies and 9% of those with tree nut allergies may outgrow them. Egg allergies affect about one in 50 children but are frequently outgrown when children reach age five. Affected individuals can be sensitive to proteins both in the egg white and egg yolk, but most children are allergic to those in the white while most adults are allergic to those in the yolk. Cow's milk is the most common food allergen in infants and young children, yet many adults are also sensitized to cow's milk. Many affected individuals cannot tolerate dairy products such as cheese and yogurt. A small portion of children with milk allergy, roughly 10%, have a reaction to beef because it contains small amounts of protein that are also present in cow's milk. Shellfish, which are divided into crustaceans (shrimp, crab, lobster, etc.) and mollusks (mussel, oyster, scallop, squid, octopus, snail, etc.), are the most common food allergy in adults. People may be allergic to other types of seafood, such as fish. Fish allergies were found to be more common in countries that have high fish consumption compared to those with lower consumption. Other common food allergens include soy and wheat. Those allergic to wheat may be sensitized to any protein in the wheat kernel. To a lesser frequency, people may be mildly allergic to raw fruits and vegetables, a disease known as oral allergy syndrome. Less common allergens include maize, spices, synthetic and natural colors, and chemical additives. Balsam of Peru, which is in various foods, is in the "top five" allergens most commonly causing patch test reactions in people referred to dermatology clinics. Routes of exposure Exposure to certain food proteins triggers the production of antigen-specific immunoglobulin E (IgE) antibodies, which, if unaccompanied by allergic symptoms, is known as allergic sensitization. Oral ingestion is the main sensitization route for most food allergy cases, yet other routes of exposure include inhalation and skin contact. For example, inhaling airborne particles in a farm-scale or factory-scale peanut shelling/crushing environment, or from cooking, can induce respiratory effects in allergic individuals. Furthermore, peanut allergies are much more common in adults who had oozing and crusted skin rashes as infants, suggesting that impaired skin may be a risk factor for sensitization. An estimated 28.5 million people worldwide are engaged in the seafood industry, which includes fishing, aquaculture, processing and industrial cooking. In these occupational settings, individuals with fish and shellfish allergies are at high risk of exposure to allergenic proteins via aerosolization. Respiratory symptoms may be induced by inhalation of wet aerosols from fresh fish handling, inhalation of dry aerosols from fishmeal processing, and dermal contact through skin breaks and cuts. Another occupational food allergy that involves respiratory symptoms is "baker's asthma," which commonly develops in food service workers who work with baked goods. Previous studies detected 40 allergens from wheat, some cross-reacted with rye proteins and a few cross-reacted with grass pollens. Allergic sensitization can occur via skin antigen exposure, which usually manifests as hives. The skin has been suggested to be a critical sensitization route for peanut-allergic individuals. Peanut allergies are much more common in adults who had oozing and crusted skin rashes as infants, reinforcing that those with disrupted epithelial barriers, notably the skin barrier, are more prone to skin sensitization. Environmental factors, such as exposure to food, microorganisms, creams, and detergents, may lead to skin barrier dysfunction. Several studies reveal that children exposed to skin creams containing peanut oil are reported to have a higher risk of peanut allergy, suggesting that impaired skin may be a risk factor for sensitization. Atopy Food allergies develop more easily in people with the atopic syndrome, a very common combination of diseases: allergic rhinitis and conjunctivitis, eczema, and asthma. The syndrome has a strong inherited component; a family history of allergic diseases can be indicative of the atopic syndrome. Cross-reactivity Some children who are allergic to cow's milk protein also show a cross-sensitivity to soy-based products. Some infant formulas have their milk and soy proteins hydrolyzed, so when taken by infants, their immune systems do not recognize the allergen and they can safely consume the product. Hypoallergenic infant formulas can be based on proteins partially predigested to a less antigenic form. Other formulas, based on free amino acids, are the least antigenic and provide complete nutritional support in severe forms of milk allergy. Crustaceans (shrimp, crab, lobster, etc.) and molluscs (mussel, oyster, scallop, squid, octopus, snail, etc.) are different invertebrate classes, but the allergenic protein tropomyosin is present and responsible for cross-reactivity. People with latex allergy often also develop allergies to bananas, kiwifruit, avocados, and some other foods. Pathophysiology Conditions caused by food allergies are classified into three groups according to the mechanism of the allergic response: IgE-mediated (classic) – the most common type, occurs shortly after eating and may involve anaphylaxis. Non-IgE mediated – characterized by an immune response not involving immunoglobulin E; may occur some hours after eating, complicating diagnosis IgE and/or non-IgE-mediated – a hybrid of the above two types Allergic reactions are abnormal immune responses to certain substances that are normally harmless. When immune cells encounter the allergenic protein, IgE antibodies are produced; this is similar to the immune system's reaction to foreign pathogens. The IgE antibodies identify the allergenic proteins as harmful and initiate the allergic reaction. The harmful proteins are those that do not break down due to the strong bonds of the protein. IgE antibodies bind to a receptor on the surface of the protein, creating a tag, just as a virus or parasite becomes tagged. Why some proteins do not denature and subsequently trigger allergic reactions and hypersensitivity while others do is not entirely clear. Hypersensitivities are categorized according to the parts of the immune system that are attacked and the amount of time it takes for the response to occur. The four types of hypersensitivity reaction are: type 1, immediate IgE-mediated; type 2, cytotoxic; type 3, immune complex-mediated; and type 4, delayed cell-mediated. The pathophysiology of allergic responses can be divided into two phases. The first is an acute response that occurs immediately after exposure to an allergen. This phase can either subside or progress into a "late-phase reaction" which can substantially prolong the symptoms of a response, and result in tissue damage. Many food allergies are caused by hypersensitivities to particular proteins in different foods. Proteins have unique properties that allow them to become allergens, such as stabilizing forces in their tertiary and quaternary structures which prevent degradation during digestion. Many theoretically allergenic proteins cannot survive the destructive environment of the digestive tract, thus do not trigger hypersensitive reactions. Acute response In the early stages of allergy, a type I hypersensitivity reaction against an allergen, encountered for the first time, causes a response in a type of immune cell called a TH2 lymphocyte, which belongs to a subset of T cells that produce a cytokine called interleukin-4 (IL-4). These TH2 cells interact with other lymphocytes called B cells, whose role is the production of antibodies. Coupled with signals provided by IL-4, this interaction stimulates the B cell to begin production of a large amount of a particular type of antibody known as IgE. Secreted IgE circulates in the blood and binds to an IgE-specific receptor (a kind of Fc receptor called FcεRI) on the surface of other kinds of immune cells called mast cells and basophils, which are both involved in the acute inflammatory response. The IgE-coated cells, at this stage, are sensitized to the allergen. If later exposure to the same allergen occurs, the allergen can bind to the IgE molecules held on the surface of the mast cells or basophils. Cross-linking of the IgE and Fc receptors occurs when more than one IgE-receptor complex interacts with the same allergenic molecule and activates the sensitized cell. Activated mast cells and basophils undergo a process called degranulation, during which they release histamine and other inflammatory chemical mediators (cytokines, interleukins, leukotrienes, and prostaglandins) from their granules into the surrounding tissue causing several systemic effects, such as vasodilation, mucous secretion, nerve stimulation, and smooth-muscle contraction. This results in rhinorrhea, itchiness, dyspnea, and anaphylaxis. Depending on the individual, the allergen, and the mode of introduction, the symptoms can be system-wide (classical anaphylaxis), or localized to particular body systems. Late-phase response After the chemical mediators of the acute response subside, late-phase responses can often occur due to the migration of other leukocytes such as neutrophils, lymphocytes, eosinophils, and macrophages to the initial site. The reaction is usually seen 2–24 hours after the original reaction. Cytokines from mast cells may also play a role in the persistence of long-term effects. Diagnosis Diagnosis is usually based on a medical history, elimination diet, skin prick test, blood tests for food-specific IgE antibodies, or oral food challenge. For skin-prick tests, a tiny board with protruding needles is used. The allergens are placed either on the board or directly on the skin. The board is then placed on the skin, to puncture the skin and for the allergens to enter the body. If a hive appears, the person is considered positive for the allergy. This test only works for IgE antibodies. Allergic reactions caused by other antibodies cannot be detected through skin-prick tests. Skin-prick testing is easy to do and results are available in minutes. Different allergists may use different devices for testing. Some use a "bifurcated needle", which looks like a fork with two prongs. Others use a "multitest", which may look like a small board with several pins sticking out of it. In these tests, a tiny amount of the suspected allergen is put onto the skin or into a testing device, and the device is placed on the skin to prick, or break through, the top layer of skin. This puts a small amount of the allergen under the skin. A hive will form at any spot where the person is allergic. This test generally yields a positive or negative result. It is good for quickly learning if a person is allergic to a particular food or not because it detects IgE. Skin tests cannot predict if a reaction would occur or what kind of reaction might occur if a person ingests that particular allergen. They can, however, confirm an allergy in light of a patient's history of reactions to a particular food. Non-IgE-mediated allergies cannot be detected by this method. Patch testing is used to determine if a specific substance causes allergic inflammation of the skin. It tests for delayed food reactions. Blood testing is another way to test for allergies; however, it poses the same disadvantage and only detects IgE allergens and does not work for every possible allergen. Radioallergosorbent testing (RAST) is used to detect IgE antibodies present to a certain allergen. The score taken from the RAST is compared to predictive values, taken from a specific type of RAST. If the score is higher than the predictive values, a great chance the allergy is present in the person exists. One advantage of this test is that it can test many allergens at one time. A CAP-RAST has greater specificity than RAST; it can show the amount of IgE present to each allergen. Researchers have been able to determine "predictive values" for certain foods, which can be compared to the RAST results. If a person's RAST score is higher than the predictive value for that food, over a 95% chance exists that patients will have an allergic reaction (limited to rash and anaphylaxis reactions) if they ingest that food. Currently, predictive values are available for milk, egg, peanut, fish, soy, and wheat. Blood tests allow for hundreds of allergens to be screened from a single sample, and cover food allergies as well as inhalants. However, non-IgE-mediated allergies cannot be detected by this method. Other widely promoted tests such as the antigen leukocyte cellular antibody test and the food allergy profile are considered unproven methods, the use of which is not advised. Food challenges test for allergens other than those caused by IgE allergens. The allergen is given to the person in the form of a pill, so the person can ingest the allergen directly. The person is watched for signs and symptoms. The problem with food challenges is that they must be performed in the hospital under careful watch, due to the possibility of anaphylaxis. Food challenges, especially double-blind, placebo-controlled food challenges, are the gold standard for diagnosis of food allergies, including most non-IgE-mediated reactions, but is rarely done. Blind food challenges involve packaging the suspected allergen into a capsule, giving it to the patient, and observing the patient for signs or symptoms of an allergic reaction. The recommended method for diagnosing food allergy is to be assessed by an allergist. The allergist will review the patient's history and the symptoms or reactions that have been noted after food ingestion. If the allergist feels the symptoms or reactions are consistent with food allergy, he/she will perform allergy tests. Additional diagnostic tools for evaluation of eosinophilic or non-IgE mediated reactions include endoscopy, colonoscopy, and biopsy. Differential diagnosis Important differential diagnoses are: Lactose intolerance generally develops later in life, but can present in young patients in severe cases. It is not an immune reaction and is due to an enzyme deficiency (lactase). It is more common in many non-Western people. Celiac disease. While it is caused by a permanent intolerance to gluten (present in wheat, rye, barley and oats), is not an allergy nor simply an intolerance, but a chronic, multiple-organ autoimmune disorder primarily affecting the small intestine. Irritable bowel syndrome C1 Esterase inhibitor deficiency (hereditary angioedema), a rare disease, generally causes attacks of angioedema, but can present solely with abdominal pain and occasional diarrhea, and thus may be confused with allergy-triggered angioedema. Prevention Breastfeeding for more than four months may prevent atopic dermatitis, cow's milk allergy, and wheezing in early childhood. Early exposure to potential allergens may be protective. Specifically, early exposure to eggs and peanuts reduces the risk of allergies to these. Guidelines suggest introducing peanuts as early as 4–6 months and include precautionary measures for high-risk infants. The former guidelines, advising delaying the introduction of peanuts, are now thought to have contributed to the increase in peanut allergy seen recently. To avoid an allergic reaction, a strict diet can be followed. It is difficult to determine the amount of allergenic food required to elicit a reaction, so complete avoidance should be attempted. In some cases, hypersensitive reactions can be triggered by exposures to allergens through skin contact, inhalation, kissing, participation in sports, blood transfusions, cosmetics, and alcohol. Early introduction of peanut and egg alongside other solids, or by one year of age, may help prevent development of food allergy. Introduction of these allergenic foods within the first year of life appears to be safe. A window of opportunity for the introduction of different food allergens may exist, such as egg introduction ahead of peanut. Inhalation exposure Allergic reactions to airborne particles or vapors of known food allergens have been reported as an occupational consequence of people working in the food industry, but can also take place in home situations, restaurants, or confined spaces such as airplanes. According to two reviews, respiratory symptoms are common, but in some cases there has been progression to anaphylaxis. The most frequent reported cases of reactions by inhalation of allergenic foods were due to peanut, seafood, legumes, tree nut, and cow's milk. Steam rising from cooking of lentils, green beans, chickpeas and fish has been well documented as triggering reactions, including anaphylactic reactions. One review mentioned case study examples of allergic responses to inhalation of other foods, including examples in which oral consumption of the food is tolerated. Treatment The mainstay of treatment for food allergy is total avoidance of the foods identified as allergens. An allergen can enter the body by consuming a portion of food containing the allergen, and can also be ingested by touching any surfaces that may have come into contact with the allergen, then touching the eyes or nose. For people who are extremely sensitive, avoidance includes avoiding touching or inhaling problematic food. Total avoidance is complicated because the declaration of the presence of trace amounts of allergens in foods is not mandatory (see regulation of labelling). If the food is accidentally ingested and a systemic reaction (anaphylaxis) occurs, then epinephrine should be used. A second dose of epinephrine may be required for severe reactions. The person should then be transported to the emergency room, where additional treatment can be given. Other treatments include antihistamines and steroids. Epinephrine Epinephrine (adrenaline) is the first-line treatment for severe allergic reactions (anaphylaxis). If administered in a timely manner, epinephrine can reverse its effects. Epinephrine relieves airway swelling and obstruction, and improves blood circulation; blood vessels are tightened and heart rate is increased, improving circulation to body organs. Epinephrine is available by prescription in an autoinjector. Antihistamines Antihistamines can alleviate some of the milder symptoms of an allergic reaction, but do not treat all symptoms of anaphylaxis. Antihistamines block the action of histamine, which causes blood vessels to dilate and become leaky to plasma proteins. Histamine also causes itchiness by acting on sensory nerve terminals. The most common antihistamine given for food allergies is diphenhydramine. Steroids Glucocorticoid steroids are used to calm down the immune system cells that are attacked by the chemicals released during an allergic reaction. This treatment in the form of a nasal spray should not be used to treat anaphylaxis, for it only relieves symptoms in the area in which the steroid is in contact. Another reason steroids should not be used is the delay in reducing inflammation. Steroids can also be taken orally or through injection, by which every part of the body can be reached and treated, but a long time is usually needed for these to take effect. Immunotherapy Immunotherapies seek to condition the immune system to elicit or suppress a specific immune response. In the treatment of allergies, common immunotherapies seek to desensitize the immune system by gradually exposing the body to allergens in increasing amounts. These forms of immunotherapy have had varying and limited success and have generally been used to treat peanut and environmental allergies. Omalizumab Omalizumab, an injectable asthma treatment drug sold under the brand name Xolair, was approved in the United States in February 2024 to reduce severe reactions to accidental exposure to food allergens. It is a genetically engineered monoclonal antibody which specifically binds to immunoglobulin E (IgE) to reduce the severity of an immune response. Successful results were reported for wheat, eggs, milk and baked products containing wheat and milk. Epidemiology Food allergies affect up to 10% of the worldwide population, and they are currently more prevalent in children (~8%) than adults (~5) in western nations. In several industrialized countries, food allergies affect up to 10% of children. Children are most commonly allergic to cow's milk, chicken eggs, peanuts, and tree nuts. While studies on adults with food allergy are not as abundant, surveys suggest that the most common food allergens for adults include fish, shellfish, peanuts, and tree nuts. Food allergies have become increasingly prevalent in industrialized/westernized nations over the last 2–3 decades. An estimated 15 million people currently have food allergies in the United States. In 1997, 0.4% children in the United States were reported to have peanut allergy, yet this number markedly rose to 1.4% in 2008. In Australia, hospital admission rates for food-induced anaphylaxis increased by an average of 13.2% from 1994-2005. One possible explanation for the rise in food allergy is the "old friends" hypothesis, which suggests that non-disease-causing organisms, such as helminths, could protect against allergy. Therefore, reduced exposure to these organisms, particularly in developed countries, could have contributed towards the increase. Children of East Asian or African descent who live in westernized countries were reported to be at significantly higher risk of food allergy compared to Caucasian children. Several studies predict that Asia and Africa may experience a growth in food allergy prevalence as the lifestyles there become more westernized. The prevalence of certain food allergies is suggested to depend partly on the geographical area and country. For instance, allergy to buckwheat flour, used for soba noodles, is more common in Japan than peanuts, tree nuts or foods made from soy beans. Also, shellfish allergy is the most common cause of anaphylaxis in adults and adolescents particularly in East Asian countries like Hong Kong, Taiwan, Singapore, and Thailand. Individuals in East Asia have further developed an allergy to rice, which forms a large part of their diet. Another example is that, out of nine European countries, egg allergy was found to be most prevalent in the UK and least prevalent in Greece. Special population: children About 75% of children who have allergies to milk protein are able to tolerate baked-in milk products, i.e., muffins, cookies, cake, and hydrolyzed formulas. About 50% of children with allergies to milk, egg, soy, peanuts, tree nuts, and wheat will outgrow their allergy by the age of 6. Those who are still allergic by the age of 12 or so have less than an 8% chance of outgrowing the allergy. United States In the United States, food allergy affects as many as 5% of infants less than three years of age and 3% to 4% of adults. The prevalence of food allergies is rising. Food allergies cause roughly 30,000 emergency room visits and 150 deaths per year. Regulation Whether rates of food allergy are increasing or not, food allergy awareness has definitely increased, with impacts on the quality of life for children, their parents and their caregivers. In the United States, the Food Allergen Labeling and Consumer Protection Act of 2004 causes people to be reminded of allergy problems every time they handle a food package, and restaurants have added allergen warnings to menus. The Culinary Institute of America, a premier school for chef training, has courses in allergen-free cooking and a separate teaching kitchen. School systems have protocols about what foods can be brought into the school. Despite all these precautions, people with serious allergies are aware that accidental exposure can easily occur at other peoples' houses, at school or in restaurants. Regulation of labelling In response to the risk that certain foods pose to those with food allergies, some countries have responded by instituting labeling laws that require food products to clearly inform consumers if their products contain priority allergens or byproducts of major allergens among the ingredients intentionally added to foods. The priority allergens vary by country. There are no labeling laws mandating declaration of the presence of trace amounts in the final product as a consequence of cross-contamination, except in Brazil. Ingredients intentionally added In the United States, the Food Allergen Labeling and Consumer Protection Act of 2004 requires companies to disclose on the label whether a packaged food product contains any of these eight major food allergens, added intentionally: cow's milk, peanuts, eggs, shellfish, fish, tree nuts, soy and wheat. The eight-ingredient list originated in 1999 from the World Health Organisation Codex Alimentarius Commission. To meet labeling requirements, if an ingredient is derived from one of the required-label allergens, then it must either have its "food sourced name" in parentheses, for example, "Casein (milk)," or as an alternative, there must be a statement separate but adjacent to the ingredients list: "Contains milk" (and any other of the allergens with mandatory labeling). The European Union requires listing for those eight major allergens plus molluscs, celery, mustard, lupin, sesame and sulfites. In 2018, the US FDA issued a request for information for the consideration of labeling for sesame to help protect people who have sesame allergies. A decision was reached in November 2020 that food manufacturers voluntarily declare that when powdered sesame seeds are used as a previously unspecified spice or flavor, the label be changed to "spice (sesame)" or "flavor (sesame)." Congress and the President passed a law in April 2021, the "FASTER Act", stipulating that labeling be mandatory, to be effect January 1, 2023, making it the ninth required food ingredient label. The Food Allergen Labeling and Consumer Protection Act of 2004 applies to packaged foods regulated by the FDA, which does not include poultry, most meats, certain egg products, and most alcoholic beverages. However, some meat, poultry, and egg processed products may contain allergenic ingredients. These products are regulated by the Food Safety and Inspection Service, which requires that any ingredient be declared in the labeling only by its common or usual name. Neither the identification of the source of a specific ingredient in a parenthetical statement nor the use of statements to alert for the presence of specific ingredients, like "Contains: milk", are mandatory. The act also does not apply to food prepared in restaurants. The EU Food Information for Consumers Regulation 1169/2011 – requires food businesses to provide allergy information on food sold unpackaged, for example, in catering outlets, deli counters, bakeries and sandwich bars. In the United States, there is no federal mandate to address the presence of allergens in drug products, medicines, or cosmetics. Trace amounts as a result of cross-contamination The value of allergen labeling other than for intentional ingredients is controversial. This concerns labeling for ingredients present unintentionally as a consequence of cross-contact or cross-contamination at any point along the food chain (during raw material transportation, storage or handling, due to shared equipment for processing and packaging, etc.). Experts in this field propose that if allergen labeling is to be useful to consumers, and healthcare professionals who advise and treat those consumers, ideally there should be agreement on which foods require labeling, threshold quantities below which labeling may be of no purpose, and validation of allergen detection methods to test and potentially recall foods that were deliberately or inadvertently contaminated. Labeling regulations have been modified to provide for mandatory labeling of ingredients plus voluntary labeling, termed precautionary allergen labeling, also known as "may contain" statements, for possible, inadvertent, trace amount, cross-contamination during production. Precautionary allergen labeling can be confusing to consumers, especially as there can be many variations on the wording of the warning. Precautionary allergen labeling is optional in the United States. , precautionary allergen labeling is regulated only in Switzerland, Japan, Argentina, and South Africa. Argentina decided to prohibit precautionary allergen labeling since 2010 and instead puts the onus on the manufacturer to control the manufacturing process and label only those allergenic ingredients known to be in the products. South Africa does not permit the use of precautionary allergen labeling, except when manufacturers demonstrate the potential presence of allergen due to cross-contamination through a documented risk assessment and despite adherence to Good Manufacturing Practice. In Australia and New Zealand there is a recommendation that precautionary allergen labeling be replaced by guidance from VITAL 2.0 (Vital Incidental Trace Allergen Labeling). A review identified "the eliciting dose for an allergic reaction in 1% of the population" as the threshold reference dose for certain foods (such as cow's milk, egg, peanut and other proteins) to provide food manufacturers with guidance for developing precautionary labeling and give consumers a better idea of what might be accidentally in a food product beyond "may contain." VITAL 2.0 was developed by the Allergen Bureau, a food industry sponsored, non-government organization. The European Union has initiated a process to create labeling regulations for unintentional contamination but is not expected to publish such before 2024. In Brazil, since April 2016, the declaration of the possibility of cross-contamination is mandatory when the product does not intentionally add any allergenic food or its derivatives, but the Good Manufacturing Practices and allergen control measures adopted are not sufficient to prevent the presence of accidental trace amounts. These allergens include wheat, rye, barley, oats and their hybrids, crustaceans, eggs, fish, peanuts, soybean, milk of all species of mammalians, almonds, hazelnuts, cashew nuts, Brazil nuts, macadamia nuts, walnuts, pecan nuts, pistachios, pine nuts, and chestnuts. Genetically modified food There is a scientific consensus that available food derived from genetically modified crops poses no greater risk to human health than conventional food, and a 2016 U.S. National Academy of Sciences report concluded that there is no relationship between consumption of genetically modified foods and the increase in prevalence of food allergies. However, there are concerns that genetically modified foods, also described as foods sourced from genetically modified organisms, could be responsible for allergic reactions, and that the widespread acceptance of these types of foods may be responsible for what is a real or perceived increase in the percentage of people with allergies. One concern is that genetic engineering could make an allergy-provoking food more allergic, meaning that smaller portions would suffice to set off a reaction. Of the food currently in widespread use of genetically modified organisms, only soybeans are identified as a common allergen. However, for the soybean proteins known to trigger allergic reactions, there is more variation from strain to strain than between those and the genetically modified varieties. Another concern is that genes transferred from one species to another could introduce an allergen in a food not thought of as particularly allergenic. Research on an attempt to enhance the quality of soybean protein by adding genes from Brazil nuts was terminated when human volunteers known to have tree nut allergy reacted to the modified soybeans. Prior to a new genetically modified food receiving government approval, certain criteria need to be met. These include: Is the donor species known to be allergenic? Does the amino acid sequence of the transferred proteins resemble the sequence of known allergenic proteins? Are the transferred proteins resistant to digestion – a trait shared by many allergenic proteins? Genes approved for animal use can be restricted from human consumption due to potential for allergic reactions. In 1998 StarLink brand corn restricted to animals was detected in the human food supply, leading to first a voluntary and then an FDA-mandated recall. There are requirements in some countries and recommendations in others that all foods containing genetically modified ingredients be so labeled, and that there be a post-launch monitoring system to report adverse effects (similar to the requirements in some countries for drug and dietary supplement reporting). Restaurants In the US, the FDA Food Code states that the person in charge in restaurants should have knowledge about major food allergens, cross-contacts, and symptoms of food allergy reactions. Restaurant staff, including wait staff and kitchen staff, may not be adequately informed about allergenic ingredients, or the risk of cross-contact when kitchen utensils used to prepare food may have been in previous contact with an allergenic food. The problem may be compounded when customers have a hard time describing their food allergies or when wait staff have a hard time understanding those with food allergies when taking an order. Diagnosing issues There exists both over-reporting and under-reporting of the prevalence of food allergies. Self-diagnosed perceptions of food allergy are greater than the rates of true food allergy because people confuse non-allergic intolerance with allergy, and also attribute non-allergy symptoms to an allergic response. Conversely, healthcare professionals treating allergic reactions on an out-patient or even hospitalized basis may not report all cases. Recent increases in reported cases may reflect a real change in incidence or an increased awareness on the part of healthcare professionals. Social impact Food fear has a significant impact on quality of life. For children with allergies, their quality of life is also affected by the actions of their peers. An increased occurrence of bullying has been observed, which can include threats or deliberate acts of forcing allergic children to contact foods that they must avoid or intentional contamination of allergen-free food. The social impacts of food allergies can carry over into adulthood. Portrayal in media Media portrayals of food allergy in television and film are not accurate, often used for comedic effect or underplaying the potential severity of an allergic reaction. These tropes misinform the public and also contribute to how entertainment media will continue to wrongly portray food allergies in the future. Types of tropes: 1) characters have food allergies, providing a weakness that can be used to sabotage them. In the movie Parasite a housekeeper is displaced by taking advantage of her peach allergy. In the animated film Peter Rabbit, the farm owner is attacked by being pelted with blackberries, causing an anaphylactic reaction requiring emergency treatment with epinephrine. After many public protests, Sony Pictures and the Peter Rabbit director apologized for making light of food allergies. 2) Food allergy is used for comedic effect, such as in the movies Hitch and in television, Kelso's egg allergy in That '70s Show. 3) Food allergies may be incorporated into characters who are also portrayed as annoying, weak and oversensitive, which can be taken as implying that their allergies are either not real or not potentially severe. In season 1, episode 16 of The Big Bang Theory Howard Wolowitz deliberately consumes a peanut-containing food bar (and has a serious reaction) just to delay Leonard from returning to his apartment where a surprise birthday party is being arranged. 4) Any of these portrayals may underplay the potential severity of food allergy, some showing that Benadryl treatment is sufficient. Viewing of humorous portrayals of food allergies has been shown to have a negative effect on related health policy support due to low perceived seriousness. Research Several theories have been suggested to explain why certain individuals develop allergic sensitization instead of oral tolerance to food allergens. One such theory is the dual allergen hypothesis, which states that ingesting food allergens early on promotes oral tolerance while skin exposure leads to sensitization. Instead of oral ingestion, skin barrier disruption in conditions like eczema, for instance, was suggested to cause allergic sensitization in animal and human infants. Inhalation was recently proposed to be an additional sensitization route in the dual allergen hypothesis. Another theory is the barrier regulation hypothesis, describing the role of commensal bacteria in preventing the development of food allergy by maintaining integrity of the intestinal epithelial barrier. Environmental and lifestyle factors, such as early life nutrition and antibiotic treatment, may contribute to food allergy prevalence by affecting gut microbial composition, and thus, intestinal immune homeostasis in infants and young children. A number of desensitization techniques are being studied. Areas of research include specific oral tolerance induction (also known as oral immunotherapy), and sublingual immunotherapy. The benefits of allergen immunotherapy for food allergies is unclear, thus is not recommended . There is research on the effects of increasing intake of polyunsaturated fatty acids during pregnancy, lactation, via infant formula and in early childhood on the subsequent risk of developing food allergies during infancy and childhood. From two reviews, maternal intake of omega-3, long-chain fatty acids during pregnancy appeared to reduce the risks of medically diagnosed IgE-mediated allergy, eczema and food allergy per parental reporting in the first 12 months of life, but the effects were not all sustained past 12 months. The reviews characterized the literature's evidence as inconsistent and limited. Results when breastfeeding mothers were consuming a diet high in polyunsaturated fatty acids were inconclusive. For infants, supplementing their diet with oils high in polyunsaturated fatty acids did not affect the risks of food allergies, eczema or asthma either as infants or into childhood. There is research on probiotics, prebiotics and the combination of the two (synbiotics) as a means of treating or preventing infant and child allergies. From reviews, there appears to be a treatment benefit for eczema, but not asthma, wheezing or rhinoconjunctivitis. The evidence was not consistent for preventing food allergies and this approach cannot yet be recommended. The Food Standards Agency, in the United Kingdom, are in charge of funding research into food allergies and intolerance. Since their founding in 1994 they have funded over 45 studies. In 2005 Europe created EuroPrevall, a multi-country project dedicated to research involving allergies.	Biology and health sciences	Specific diseases	Health
679597	https://en.wikipedia.org/wiki/Vajont%20Dam	Vajont Dam	The Vajont Dam or Vaiont Dam is a disused hydro-electric dam in northern Italy. It is one of the tallest dams in the world, with a height of . It is in the valley of the Vajont (river) under Monte Toc, in the municipality of Erto e Casso, north of Venice. The dam was conceived in the 1920s and eventually built between 1957 and 1960 by Società Adriatica di Elettricità, at the time the electricity supply and distribution monopoly in northeastern Italy. The engineer was Carlo Semenza (1893–1961). In 1962, the dam was nationalized and came under the control of ENEL as part of the Italian Ministry of Public Works. On 9 October 1963, during initial filling of the lake, a landslide caused a megatsunami in which of water overtopped the dam in a wave of , bringing massive flooding and destruction to the Piave Valley below, destroying several villages and towns, causing an estimated 1,900 to 2,500 deaths. The dam itself remained almost intact and two-thirds of the water was retained behind it. This event occurred after ENEL and the Italian government concealed reports and dismissed evidence that Monte Toc, on the southern side of the lake, was geologically unstable. They had disregarded numerous warnings, danger signals, and negative appraisals. Underestimating the size of the landslide, ENEL's attempt to safely mitigate any landslide by lowering the level of the lake came too late, when disaster was almost imminent. Construction The dam was built by Società Adriatica di Elettricità (SADE), the electricity supply and distribution monopoly in northeastern Italy. The owner, Giuseppe Volpi di Misurata, had been Mussolini's minister of finance for several years. The "tallest dam in the world", across the Vajont Gorge, was conceived in the 1920s to harness the Piave, Mae, and Boite Rivers, to meet the growing demand for power generation and industrialization. The project was not authorized until 15 October 1943, after the confusion following Mussolini's fall during World War II was resolved. The dam wall had a volume of and held up to of water. The dam and basin were intended to be at the centre of a complex water-management system in which water would have been channeled from nearby valleys and artificial basins at higher levels. Tens of kilometres of concrete pipes and pipe bridges across valleys were planned. In the 1950s, SADE's monopoly was confirmed by post-fascist governments, and it bought the land despite opposition from the communities of Erto and Casso in the valley, which was overcome by government and police support. SADE stated that the geology of the gorge had been studied, including analysis of previous landslides, and that the mountain was held to be sufficiently stable. Construction work on the dam started in 1957. By 1959, shifts and fractures in the ground were noticed while building a new road on the side of Monte Toc. This led to further studies in which three experts separately told SADE that the entire side of Monte Toc was unstable and was likely to collapse into the basin if the filling were completed, due to the raised water level increasing the instability. SADE took no notice and construction was completed in October 1959. In February 1960, SADE was authorised to start filling the lake. In 1962, the dam was nationalized and came under the control of ENEL as part of the Italian Ministry for Public Works. Early signs of disaster On 22 March 1959, during construction of the Vajont Dam, a landslide at the nearby Pontesei Dam created a wave that killed one person. Throughout the summer of 1960, minor landslides and earth movements were noticed. Instead of heeding these warning signs, the Italian government chose to sue the handful of journalists reporting the problems for "undermining the social order". On 4 November 1960, with the water level in the reservoir at about of the planned , a landslide of about collapsed into the lake. SADE stopped the filling, lowered the water level by about , and started to build an artificial gallery in the basin in front of Monte Toc to keep the basin usable even if additional landslides (which were expected) divided it into two parts. In October 1961, after the completion of the gallery, SADE resumed filling the reservoir under controlled monitoring. In April and May 1962, with the basin water level at , the people of Erto and Casso reported five "grade five" Mercalli intensity scale earthquakes. SADE downplayed the importance of these quakes, and was then authorized to fill the reservoir to the maximum level. In July 1962, SADE's own engineers reported the results of model-based experiments on the effects of further landslides from Monte Toc into the lake. The tests indicated that a wave generated by a landslide could overtop the crest of the dam if the water level was or less from the dam crest. SADE therefore decided that a reduced level below the crest would prevent any displacement wave from over-topping the dam. In March 1963, the dam was transferred to the newly constituted government service for electricity, ENEL. During the following summer, with the reservoir almost completely filled, landslips, shakes, and movements of the ground were repeatedly reported by the alarmed population. On 15 September 1963, the entire side of the mountain slid by . On 26 September 1963, ENEL decided to slowly empty the basin to . By early October 1963, the collapse of the mountain's south side looked unavoidable; in one day, it moved almost . Tests of the hydraulic model of the Vajont tank After the discovery of the landslide on the northern slopes of Monte Toc, it was decided to deepen the studies on these effects: Dynamic actions on the dam Wave effects in the reservoir and possible dangers for nearby locations, with particular attention to the town of Erto Hypothesis of a partial breakage of the dam and consequent examination of the rout wave and its propagation along the last stretch of the Vajont and along the Piave, up to Soverzene and beyond The study of point 1 was performed at the Experimental Institute for Models and Structures (ISMES) of Bergamo, while for the others, the SADE decided to build a physical-hydraulic model of the basin, in which to perform some experiments on the effects of a landslide fall in a reservoir. The 1: 200 scale model of the basin, which can still be visited today, was set up at the SADE hydroelectric plant in Nove (Borgo Botteon di Vittorio Veneto), and became the Hydraulic Models Centre. The experiments were entrusted to professors Ghetti and Marzolo, university professors of the Institute of Hydraulics and Hydraulic Constructions of the University of Padua, and were carried out with funding of SADE, under the control of the study office of the company itself. The study aimed to verify the hydraulic effects on the dam and on the banks of the landslide reservoir, so was directed in this sense rather than reproducing the natural phenomenon of the landslide. The experiments were carried out in two different series (August–September 1961 and January–April 1962), of which the first served substantially to refine the model. First set of experiments The first series of five experiments began on 30 August 1961, with a sliding surface of the flat landslide inclined by 30°, consisting of a wooden plank covered with a sheet. The sliding mass was simulated with gravel, held in place by flexible metal nets, which were initially held in position by ropes that were then suddenly released. At the beginning of September, another four tests were carried out intended for orientation purposes. The first always with a 30° inclined plane, the following three with a 42° inclined plane. Having found it impossible to reproduce the natural geological phenomenon of the landslide in the model, the model was elaborated by modifying the movement surface of the landslide, which was replaced with a masonry one (the relative profiles were elaborated by Semenza, who also used the surveys that had already been carried out and that had provided sufficient elements of judgment in this sense), to make it possible to vary the speed of the landslide fall into the reservoir (made difficult by the new "back" shape of the movement surface). To simulate the compactness of the moving material (which in the model remained the gravel), rigid sectors were inserted that were towed by ropes pulled by a tractor. Second set of experiments In the second set of 17 experiments, conducted from 3 January 1962, to 24 April 1962, the "collapsing" material was still gravel, this time held in place by hemp nets and cords. Starting from the Muller hypothesis relating to the different characteristics of the mass moving between the downstream part of the Massalezza stream (west) and the upstream part of the same (east), all the experiments were performed by making those two hypothetical parts of the landslide descend separately. In the model, however, the two landslides were initially made to descend at different times, so that their effects were totally separate, and subsequently, when the wave produced by the first came back, so as to obtain a total increase in the water of the even greater lake. Final Ghetti report The total increase of the water in the tank (measured by means of special instruments) was broken down into "static increase", which was the nontransient effect of increasing the level of the water left in the tank after the landslide due to the immersion of the landslide in the tank (once the state of rest is reached again), and in "dynamic boost", due to the temporary wave motion produced by the landslide. The static increase depended on the volume of the landslide that remained immersed in the tank, while the dynamic boost depended almost exclusively on the speed of the landslide fall (while it was negligibly linked to the volume of the same). Based on this simulation (following the disaster object of criticism, as considered approximate by some), placing a reservoir limit at an altitude of predicted that no damage would occur above above mean sea level along the banks of the reservoir, while a minimum quantity of water would have exceeded the edge of the dam (), causing negligible damage downstream of the same. With the reported experiences, carried out on a 1:200 scale model of the Vajont lake-reservoir, we tried to provide an evaluation of the effects that will be caused by a landslide, which is possible to occur on the left bank upstream of the dam. . Given that the extreme limit downstream of the landslide is more than 75 m from the embankment of the dam, and that the formation of this embankment is of compact and consistent rock and well distinct, even geologically, from the aforementioned mass, it is absolutely not to be fear of any static perturbation to the dam with the occurrence of the landslide, and therefore only the effects of the wave rise in the lake and in the overflow on the dam crest as a consequence of the fall are to be considered. Landslide and wave The alarming rate of movement of the landslide had not slowed as a result of lowering the water level, although the water had been lowered to what SADE believed was a safe level to contain the displacement wave should a catastrophic landslide occur. On 9 October 1963, engineers witnessed trees falling and rocks rolling down into the reservoir where the predicted landslide would take place. With a major landslide now seeming imminent, engineers gathered at the site that evening to witness the tsunami which, when it happened, killed 60 workers. At 10:39 pm, a massive landslide, with around of trees, earth, and rock, collapsed from the northern flank of Monte Toc into the lake below at up to (another source gives ), generating a seismic shock. In 20 seconds it reached the water level; by 45 seconds the landslide (now at rest) had completely filled the Vajont reservoir. The impact displaced of water in about 25 seconds, of which overtopped the dam in a wave. Impact of the landslip with the water generated three waves. One went upwards, reached the houses of Casso, fell back onto the landslide, and travelled onwards to dig the basin of the pond of Massalezza. Another wave headed toward the shores of the lake, and by washout action, destroyed some localities in the municipality of Erto e Casso. The third wave (containing about of water) rose over the edge of the dam, which remained intact except for the roadway that led to the left side of the Vajont, and fell into the narrow valley below the dam. The roughly of water that climbed over the dam reached the stony shore of the Piave Valley and swept up substantial debris, which poured into the southern sector of Longarone and destroyed the town except for the town hall, the houses north of it, and other neighboring towns. The death toll was about 2,000 people (official data speak of 1,917 victims, but determining an accurate number with certainty is impossible). Firefighters who set out from Belluno, after reports of the raising of the level of the Piave, could not reach the location, since from a certain point onwards, the road coming from the valley had been completely swept away. Longarone was reached by firefighters who departed Pieve di Cadore, who were the first to realize what had happened and were able to communicate it. At 5:30 am on 10 October 1963, the first soldiers of the Italian Army arrived to bring relief and recover the dead. The soldiers involved were mostly Alpini, some of whom belonged to the combat engineers, who dug by hand to seek the bodies of the missing. They also found safes of the banks of the country, no longer able to be opened with normal keys, as they were damaged. The firefighters from 46 Provincial Commands also participated in the rescue, with 850 men, including divers, land and helicopter teams, and many vehicles and equipment. The Nucleo Sommozzatori of Genoa, with eight personnel, were used in the basin in front of the Busche Dam, to dredge for bodies and drums of toxic substances (61 cyanide drums), with subsequent patrol by immersion and final removal of sludge when the basin was dry. Of the nearly 2,000 fatalities, only 1,500 bodies were recovered and recomposed, of which half were beyond identification. In the Piave Valley the wave destroyed the villages of Longarone, Pirago, Rivalta, Villanova, and Faè, killing approximately 2,000 people and turning the land below the dam into a flat expanse of mud with an impact crater deep and wide. Many small villages near the landslide along the lakefront suffered damage from a giant displacement wave. Villages in the territory of Erto e Casso and the village of , near Castellavazzo, were largely destroyed. Estimates of the dead range from the official estimate of 1,917 to 2,500 people. About 350 families lost all family members. Most of the survivors had lost relatives and friends along with their homes and belongings. The dam itself was largely undamaged. The top or so of masonry was washed away, but the basic structure remained intact and is standing today. Causes and responsibility Immediately after the disaster, the government (which owned the dam), politicians, and public authorities insisted on attributing the tragedy to an unexpected and unavoidable natural event. The debate in the newspapers was heavily influenced by politics. The paper l'Unità, the mouthpiece of the Partito Comunista Italiano (PCI), was the first to denounce the actions of the management and government, as it had previously carried a number of articles by addressing the behaviour of the SADE management in the Vajont project and elsewhere. Indro Montanelli, then the most influential Italian journalist and a vocal anticommunist, attacked l'Unità and denied any human responsibility; l'Unità and the PCI were dubbed "jackals, speculating on pain and on the dead" in many articles by La Domenica del Corriere and a national campaign poster paid for by the Christian Democracy, the party of Prime Minister Giovanni Leone. They attributed the catastrophe only to natural causes and God's will. The campaign accused the PCI of sending agitprops into the refugee communities, as relief personnel; most of them were partisans from Emilia Romagna who fought on Mount Toc in the Second World War and often had friends in the stricken area. The Christian Democracy accused the Communist Party of "political profiteering" from the tragedy. Leone promised to bring justice to the people killed in the disaster. A few months after he lost the premiership, he became the head of SADE's team of lawyers, who significantly reduced the amount of compensation for the survivors and ruled out payment for at least 600 victims. The DC's newspaper, La Discussione, called the disaster "a mysterious act of God's love", in an article that drew sharp criticism from l'Unità. Apart from journalistic attacks and the attempted cover-up from news sources aligned with the government, flaws in the geological assessments have been proven, with disregard of warnings about the likelihood of a disaster by SADE, ENEL, and the government. The trial was moved to L'Aquila, in Abruzzo, by the judges who heard the preliminary trial, thus preventing public participation, and resulted in lenient sentencing for a handful of the SADE and ENEL engineers. One SADE engineer (Mario Pancini) committed suicide in 1968. The government never sued SADE for damage compensation. Subsequent engineering analysis has focused on the cause of the landslide, and debate continues about the contribution of rainfall, dam level changes, and earthquakes as triggers of the landslide, as well as differing views about whether it was an old landslide that slipped further or a completely new one. A number of problems existed with the choice of site for the dam and reservoir; the canyon was steep-sided, the river had undercut its banks, and the limestone and claystone rocks that made up the walls of the canyon were interbedded with the slippery, clay-like Lias and Dogger Jurassic-period horizons and the Cretaceous-period Malm horizon, all of which were inclined towards the axis of the canyon. In addition, the limestone layers contained many solution caverns that became only more saturated because of rains in September. Prior to the landslide that caused the overtopping flood, the downhill creep of the regolith had been per week. In September, this creep reached per day until finally, the day before the landslide, the creep was measured at . Reconstruction Most of the survivors were moved into a newly built village, Vajont, southeast on the Tagliamento River plain. Those who insisted on returning to their mountain life in Erto e Casso were strongly discouraged. Longarone and other villages in the Piave Valley were rebuilt with modern houses and farms. The government used the disaster to promote the industrialization of northeastern Italy. Survivors were entitled to "business start-up" loans, public subsidies, and 10 years of tax exemption, all of which they could "sell on" to major companies from the Venice region. These concessions were then converted into millions of lira for industrial plants elsewhere. Among the corporations were Zanussi (now owned by Electrolux), Ceramica Dolomite (now owned by American Standard), Confezioni SanRemo, and SAVIC (now owned by Italcementi). Compensation measures did not clearly differentiate between victims and people who lived nearby; thus, much of the compensation went to people who had suffered little damage, creating a negative public image. A pumping station was installed in the dam basin to keep the lake at a constant level, and the bypass gallery was lengthened beyond the dam to let the water flow down to the Piave Valley. The dam wall is still in place and maintained, but no plans exist to exploit it for electric power. The reservoir behind the dam, now dry and filled with landslip, has been open to visitors since 2002. The dam today and memorials In recent years, revival of interest has occurred both by researchers with specialist interest and sightseers. The dam, now owned by ENEL, was partially opened to the public in 2002 with guided tours and access to the walkway along the top and other locations. In September 2006, an annual noncompetitive track event, called "Paths of Remembrance", was inaugurated, which allows participants to access some locations inside the mountain. On 12 February 2008, in launching the International Year of Planet Earth, UNESCO cited the Vajont Dam tragedy as one of five "cautionary tales" caused by "the failure of engineers and geologists". For 2013, on the occasion of the 50th anniversary of the disaster, the region of Venice set aside one million euros for safety works and recovery of tunnels inside the mountain, which were part of the Colomber Road (the old national road 251). The memorial church in Longarone—constructed in spite of the strong opposition of the surviving parish priest — was built by architect Giovanni Michelucci. In the media After the initial worldwide reporting, the tragedy became regarded as part of the price of economic growth in the 1950s and 1960s. Interest was rejuvenated by a 1997 television program by Marco Paolini and , Il racconto del Vajont. In 2001, a docudrama about the disaster was released. A joint production of Italian and French companies, it was titled Vajont—La diga del disonore ("Vajont—The Dam of Dishonour") in Italy, and La Folie des hommes (The Madness of Men) in France. It stars Michel Serrault and Daniel Auteuil. The tragedy was included in the 2008 documentary series Disasters of the Century. The TV show Seconds from Disaster featured the event in episode two, "Mountain Tsunami", of its fifth season in 2012. In 2013, the 11th stage of the Giro d'Italia finished in Vajont to commemorate the 50th anniversary of the disaster. In March 2018, the dam and the disaster were also covered in season two, episode one ("Armageddon Highway") of Science Channel's Mysteries of the Abandoned.	Technology	Dams	null
679709	https://en.wikipedia.org/wiki/Nonlinear%20programming	Nonlinear programming	In mathematics, nonlinear programming (NLP) is the process of solving an optimization problem where some of the constraints are not linear equalities or the objective function is not a linear function. An optimization problem is one of calculation of the extrema (maxima, minima or stationary points) of an objective function over a set of unknown real variables and conditional to the satisfaction of a system of equalities and inequalities, collectively termed constraints. It is the sub-field of mathematical optimization that deals with problems that are not linear. Definition and discussion Let n, m, and p be positive integers. Let X be a subset of Rn (usually a box-constrained one), let f, gi, and hj be real-valued functions on X for each i in {1, ..., m} and each j in {1, ..., p}, with at least one of f, gi, and hj being nonlinear. A nonlinear programming problem is an optimization problem of the form Depending on the constraint set, there are several possibilities: feasible problem is one for which there exists at least one set of values for the choice variables satisfying all the constraints. an infeasible problem is one for which no set of values for the choice variables satisfies all the constraints. That is, the constraints are mutually contradictory, and no solution exists; the feasible set is the empty set. unbounded problem is a feasible problem for which the objective function can be made to be better than any given finite value. Thus there is no optimal solution, because there is always a feasible solution that gives a better objective function value than does any given proposed solution. Most realistic applications feature feasible problems, with infeasible or unbounded problems seen as a failure of an underlying model. In some cases, infeasible problems are handled by minimizing a sum of feasibility violations. Some special cases of nonlinear programming have specialized solution methods: If the objective function is concave (maximization problem), or convex (minimization problem) and the constraint set is convex, then the program is called convex and general methods from convex optimization can be used in most cases. If the objective function is quadratic and the constraints are linear, quadratic programming techniques are used. If the objective function is a ratio of a concave and a convex function (in the maximization case) and the constraints are convex, then the problem can be transformed to a convex optimization problem using fractional programming techniques. Applicability A typical non-convex problem is that of optimizing transportation costs by selection from a set of transportation methods, one or more of which exhibit economies of scale, with various connectivities and capacity constraints. An example would be petroleum product transport given a selection or combination of pipeline, rail tanker, road tanker, river barge, or coastal tankship. Owing to economic batch size the cost functions may have discontinuities in addition to smooth changes. In experimental science, some simple data analysis (such as fitting a spectrum with a sum of peaks of known location and shape but unknown magnitude) can be done with linear methods, but in general these problems are also nonlinear. Typically, one has a theoretical model of the system under study with variable parameters in it and a model the experiment or experiments, which may also have unknown parameters. One tries to find a best fit numerically. In this case one often wants a measure of the precision of the result, as well as the best fit itself. Methods for solving a general nonlinear program Analytic methods Under differentiability and constraint qualifications, the Karush–Kuhn–Tucker (KKT) conditions provide necessary conditions for a solution to be optimal. If some of the functions are non-differentiable, subdifferential versions of Karush–Kuhn–Tucker (KKT) conditions are available. Under convexity, the KKT conditions are sufficient for a global optimum. Without convexity, these conditions are sufficient only for a local optimum. In some cases, the number of local optima is small, and one can find all of them analytically and find the one for which the objective value is smallest. Numeric methods In most realistic cases, it is very hard to solve the KKT conditions analytically, and so the problems are solved using numerical methods. These methods are iterative: they start with an initial point, and then proceed to points that are supposed to be closer to the optimal point, using some update rule. There are three kinds of update rules: Zero-order routines - use only the values of the objective function and constraint functions at the current point; First-order routines - use also the values of the gradients of these functions; Second-order routines - use also the values of the Hessians of these functions. Third-order routines (and higher) are theoretically possible, but not used in practice, due to the higher computational load and little theoretical benefit. Branch and bound Another method involves the use of branch and bound techniques, where the program is divided into subclasses to be solved with convex (minimization problem) or linear approximations that form a lower bound on the overall cost within the subdivision. With subsequent divisions, at some point an actual solution will be obtained whose cost is equal to the best lower bound obtained for any of the approximate solutions. This solution is optimal, although possibly not unique. The algorithm may also be stopped early, with the assurance that the best possible solution is within a tolerance from the best point found; such points are called ε-optimal. Terminating to ε-optimal points is typically necessary to ensure finite termination. This is especially useful for large, difficult problems and problems with uncertain costs or values where the uncertainty can be estimated with an appropriate reliability estimation. Implementations There exist numerous nonlinear programming solvers, including open source: ALGLIB (C++, C#, Java, Python API) implements several first-order and derivative-free nonlinear programming solvers NLopt (C/C++ implementation, with numerous interfaces including Julia, Python, R, MATLAB/Octave), includes various nonlinear programming solvers SciPy (de facto standard for scientific Python) has scipy.optimize solver, which includes several nonlinear programming algorithms (zero-order, first order and second order ones). IPOPT (C++ implementation, with numerous interfaces including C, Fortran, Java, AMPL, R, Python, etc.) is an interior point method solver (zero-order, and optionally first order and second order derivatives). Numerical Examples 2-dimensional example A simple problem (shown in the diagram) can be defined by the constraints with an objective function to be maximized where . 3-dimensional example Another simple problem (see diagram) can be defined by the constraints with an objective function to be maximized where .	Mathematics	Other	null
679929	https://en.wikipedia.org/wiki/Stone%20%28unit%29	Stone (unit)	The stone or stone weight (abbreviation: st.) is an English and British imperial unit of mass equal to 14 avoirdupois pounds (6.35 kg). The stone continues in customary use in the United Kingdom and Ireland for body weight. England and other Germanic-speaking countries of Northern Europe formerly used various standardised "stones" for trade, with their values ranging from about 5 to 40 local pounds (2.3 to 18.1 kg) depending on the location and objects weighed. With the advent of metrication, Europe's various "stones" were superseded by or adapted to the kilogram from the mid-19th century onward. Antiquity The name "stone" derives from the historical use of stones for weights, a practice that dates back into antiquity. The Biblical law against the carrying of "diverse weights, a large and a small" is more literally translated as "you shall not carry a stone and a stone (), a large and a small". There was no standardised "stone" in the ancient Jewish world, but in Roman times stone weights were crafted to multiples of the Roman pound. Such weights varied in quality: the Yale Medical Library holds 10- and 50-pound examples of polished serpentine, while a 40-pound example at the Eschborn Museum is made of sandstone. Great Britain and Ireland The 1772 edition of the Encyclopædia Britannica defined the stone:STONE also denotes a certain quantity or weight of some commodities. A stone of beef, in London, is the quantity of eight pounds; in Hertfordshire, twelve pounds; in Scotland sixteen pounds. The Weights and Measures Act 1824 (5 Geo. 4. c. 74), which applied to all of the United Kingdom of Great Britain and Ireland, consolidated the weights and measures legislation of several centuries into a single document. It revoked the provision that bales of wool should be made up of 20 stones, each of 14 pounds, but made no provision for the continued use of the stone. Ten years later, a stone still varied from 5 pounds (glass) to 8 pounds (meat and fish) to 14 pounds (wool and "horseman's weight"). The Weights and Measures Act 1835 permitted using a stone of 14 pounds for trade but other values remained in use. James Britten, in 1880 for example, catalogued a number of different values of the stone in various British towns and cities, ranging from 4 lb to 26 lb. The value of the stone and associated units of measure that were legalised for purposes of trade were clarified by the Weights and Measures Act 1835 as follows: England The English stone under law varied by commodity and in practice varied according to local standards. The Assize of Weights and Measures, a statute of uncertain date from , describes stones of 5 merchants' pounds used for glass; stones of 8 lb. used for beeswax, sugar, pepper, alum, cumin, almonds, cinnamon, and nutmegs; stones of 12 lb. used for lead; and the of lb. used for wool. In 1350 Edward III issued a new statute defining the stone weight, to be used for wool and "other Merchandizes", at 14 pounds, reaffirmed by Henry VII in 1495. In England, merchants traditionally sold potatoes in half-stone increments of 7 pounds. Live animals were weighed in stones of 14 lb; but, once slaughtered, their carcasses were weighed in stones of 8 lb. Thus, if the animal's carcass accounted for of the animal's weight, the butcher could return the dressed carcasses to the animal's owner stone for stone, keeping the offal, blood and hide as his due for slaughtering and dressing the animal. Smithfield market continued to use the 8 lb stone for meat until shortly before the Second World War. The Oxford English Dictionary also lists: Scotland The Scottish stone was equal to 16 Scottish pounds (17 lb 8 oz avoirdupois or 7.936 kg). In 1661, the Royal Commission of Scotland recommended that the Troy stone be used as a standard of weight and that it be kept in the custody of the burgh of Lanark. The tron (or local) stone of Edinburgh, also standardised in 1661, was 16 tron pounds (22 lb 1 oz avoirdupois or 9.996 kg). In 1789 an encyclopedic enumeration of measurements was printed for the use of "his Majesty's Sheriffs and Stewards Depute, and Justices of Peace, ... and to the Magistrates of the Royal Boroughs of Scotland" and provided a county-by-county and commodity-by-commodity breakdown of values and conversions for the stone and other measures. The Scots stone ceased to be used for trade when the Weights and Measures Act 1824 (5 Geo. 4. c. 74) established a uniform system of measure across the whole of the United Kingdom, which at that time included all of Ireland. Ireland Before the early 19th century, as in England, the stone varied both with locality and with commodity. For example, the Belfast stone for measuring flax equaled 16.75 avoirdupois pounds. The most usual value was 14 pounds. Among the oddities related to the use of the stone was the practice in County Clare of a stone of potatoes being 16 lb in the summer and 18 lb in the winter. Modern use In 1965, the Federation of British Industry informed the British government that its members favoured adopting the metric system. The Board of Trade, on behalf of the government, agreed to support a ten-year metrication programme. There would be minimal legislation, as the programme was to be voluntary and costs were to be borne where they fell. Under the guidance of the Metrication Board, the agricultural product markets achieved a voluntary switchover by 1976. The stone was not included in the Directive 80/181/EEC as a unit of measure that could be used within the EEC for "economic, public health, public safety or administrative purposes", though its use as a "supplementary unit" was permitted. The scope of the directive was extended to include all aspects of the EU internal market from 1 January 2010. With the adoption of metric units by the agricultural sector, the stone was, in practice, no longer used for trade; and, in the Weights and Measures Act 1985, passed in compliance with EU directive 80/181/EEC, the stone was removed from the list of units permitted for trade in the United Kingdom. In 1983, in response to the same directive, similar legislation was passed in Ireland. The act repealed earlier acts that defined the stone as a unit of measure for trade. (British law had previously been silent regarding other uses of the stone.) The stone remains widely used in the United Kingdom and Ireland for human body weight: in those countries people may commonly be said to weigh, e.g., "11 stone 4" (11 stones and 4 pounds), rather than "72 kilograms" as in most of the other countries, or "158 pounds", the conventional way of expressing the same weight in the US and in Canada. The invariant plural form of stone in this context is stone (as in, "11 stone" or "12 stone 6 pounds"); in other contexts, the correct plural is stones (as in, "Please enter your weight in stones and pounds"). In Australia and New Zealand, metrication has entirely displaced stones and pounds since the 1970s. In many sports in both the UK and Ireland, such as professional boxing, wrestling, and horse racing, the stone is used to express body weights. Elsewhere The use of the stone in the former British Empire was varied. In Canada for example, it never had a legal status. Shortly after the United States declared independence, Thomas Jefferson, then Secretary of State, presented a report on weights and measures to the U.S. House of Representatives. Even though all the weights and measures in use in the United States at the time were derived from English weights and measures, his report made no mention of the stone being used. He did, however, propose a decimal system of weights in which his "[decimal] pound" would have been and the "[decimal] stone" would have been . Before the advent of metrication, units called "stone" (; ; ) were used in many northwestern European countries. Its value, usually between 3 and 10 kg, varied from city to city and sometimes from commodity to commodity. The number of local "pounds" in a stone also varied from city to city. During the early 19th century, states such as the Netherlands (including Belgium) and the South Western German states, which had redefined their system of measures using the as a reference for weight (mass), also redefined their stone to align it with the kilogram. This table shows a selection of stones from various northern European cities: Metric stone In the Netherlands, where the metric system was adopted in 1817, the pond (pound) was set equal to half a kilogram, and the steen (stone), which had previously been 8 Amsterdam pond (3.953 kg), was redefined as being 3 kg. In modern colloquial Dutch, a pond is used as an alternative for 500 grams or half a kilogram, while the ons is used for a weight of 100 grams, being pond.	Physical sciences	Mass and weight	Basics and measurement
680871	https://en.wikipedia.org/wiki/Millau%20Viaduct	Millau Viaduct	The Millau Viaduct (, ) is a multispan cable-stayed bridge completed in 2004 across the gorge valley of the Tarn near (west of) Millau in the Aveyron department in the Occitanie Region, in Southern France. The design team was led by engineer Michel Virlogeux and English architect Norman Foster. it is the tallest bridge in the world, having a structural height of . The Millau Viaduct is part of the A75–A71 autoroute axis from Paris to Béziers and Montpellier. The cost of construction was approximately (). It was built over three years, formally inaugurated on 14 December 2004, and opened to traffic two days later on 16 December. The bridge has been consistently ranked as one of the greatest engineering achievements of modern times, and received the 2006 Outstanding Structure Award from the International Association for Bridge and Structural Engineering. History In the 1980s, high levels of road traffic near Millau in the Tarn valley were causing congestion, especially in the summer due to holiday traffic on the route from Paris to Spain. A method of bypassing Millau had long been considered, not only to ease the flow and reduce journey times for long-distance traffic, but also to improve the quality of access to Millau for its local businesses and residents. One of the solutions considered was the construction of a road bridge to span the river and gorge valley. The first plans for a bridge were discussed in 1987 by CETE, and by October 1991 the decision was made to build a high crossing of the Tarn by a structure of around in length. During 1993–1994, the government consulted with seven architects and eight structural engineers. During 1995–1996, a second definition study was made by five associated architect groups and structural engineers. In January 1995, the government issued a declaration of public interest to solicit design approaches for a competition. In July 1996 the jury decided in favour of a cable-stayed design with multiple spans, as proposed by the SODETEG consortium led by Michel Virlogeux, Norman Foster and Arcadis. The decision to proceed by grant of contract was made in May 1998; then in June 2000, the contest for the construction contract was launched, open to four consortia. In March 2001, Eiffage established the subsidiary Compagnie Eiffage du Viaduc de Millau (CEVM), and was declared winner of the contest and awarded the prime contract in August. Possible routes In initial studies, four potential options were examined: Great Eastern () (yellow route) – passing east of Millau and crossing the valleys of the Tarn and Dourbie on two very high and long bridges (spans of ) whose construction was acknowledged to be problematic. This option would have allowed access to Millau only from the Larzac plateau, using the long and tortuous descent from La Cavalerie. Although this option was shorter and better suited to through-traffic, it did not satisfactorily serve the needs of Millau and its area. Great Western () (black route) – longer than the eastern option by , following the Cernon valley. Technically easier (requiring four viaducts), this solution was judged to have negative impacts on the environment, in particular on the picturesque villages of Peyre and Saint-Georges-de-Luzençon. It was more expensive than the preceding option, and served the region badly. Near RN9 () (red route) – would have served the town of Millau well, but presented technical difficulties, and would have had a strong impact on existing or planned structures. Intermediate (), west of Millau (blue route) – was supported by local opinion, but presented geological difficulties, notably on the question of crossing the valley of the Tarn. Expert investigation concluded that these obstacles were not insurmountable. The fourth option was selected by ministerial decree on 28 June 1989. It encompassed two possibilities: the high solution, envisaging a viaduct more than above the river; the low solution, descending into the valley and crossing the river on a bridge, then a viaduct of , extended by a tunnel on the Larzac side. After long construction studies by the Ministry of Public Works, the low solution was abandoned because it would have intersected the water table, had a negative impact on the town, cost more, and lengthened the driving distance. The choice of the 'high' solution was decided by ministerial decree on 29 October 1991. After the choice of the high viaduct, five teams of architects and researchers worked on a technical solution. The concept and design for the bridge was devised by French designer and structural engineer Michel Virlogeux. He worked with the Dutch engineering firm Arcadis, responsible for the structural engineering of the bridge. Choosing the definitive route The 'high solution' required the construction of a viaduct. From 1991 to 1993, the structures division of Sétra, directed by Virlogeux, carried out preliminary studies, and examined the feasibility of a single structure spanning the valley. Taking into account technical, architectural, and financial issues, the Administration of Roads opened the question for competition among structural engineers and architects to widen the search for realistic designs. By July 1993, seventeen structural engineers and thirty-eight architects applied as candidates for the preliminary studies. With the assistance of a multidisciplinary commission, the Administration of Roads selected eight structural engineers for a technical study, and seven architects for the architectural study. Choice of technical design Simultaneously, a school of international experts representing a wide spectrum of expertise (technical, architectural, and landscape), chaired by Jean-François Coste, was established to clarify the choices that had to be made. In February 1995, on the basis of proposals of the architects and structural engineers, and with support of the school of experts, five general designs were identified. The competition was relaunched: five combinations of architects and structural engineers, drawn from the best candidates of the first phase, were formed; each was to conduct in-depth studies of one of the general designs. On 15 July 1996, Bernard Pons, minister of Public Works, announced the decision of the jury, which was constituted of elected artists and experts, and chaired by Christian Leyrit, the director of highways. The solution of a multiple-span viaduct cable-stayed bridge, presented by the structural engineering group Sogelerg, Europe Etudes Gecti and Serf, and the architects Foster + Partners was declared the best. Detailed studies were carried out by the successful consortium, steered by the highways authority until mid-1998. After undergoing wind tunnel tests, the shape of the road deck was altered, and detailed corrections were made to the design of the pylons. When the details were eventually finalised, the whole design was approved in late 1998. Contractors Once the Ministry of Public Works had taken the decision to offer the construction and operation of the viaduct as a grant of contract, an international call for tenders was issued in 1999. Five consortia tendered: Compagnie Eiffage du Viaduc de Millau (CEVM), a new subsidiary created by Eiffage; PAECH Construction Enterprise, Poland; a consortium led by the Spanish company Dragados, with Skanska, Sweden, and Bec, France; Société du Viaduc de Millau, including the French companies ASF, Egis Projects, GTM Construction, Bouygues Travaux Publics, SGE, CDC Projets, Tofinso, and the Italian company Autostrade; a consortium led by Générale Routière, with Via GTI (France) and Cintra, Nesco, Acciona, and Ferrovial Agroman (Spain). Piers were built with Lafarge high performance concrete. The pylons of the Millau Viaduct, which are the tallest elements (the tallest one being ) were produced and mounted by PAECH Construction Enterprise from Poland. The Compagnie Eiffage du Viaduc de Millau, working with the architect Norman Foster, was successful in obtaining the tender. Because the government had already taken the design work to an advanced stage, the technical uncertainties were considerably reduced. A further advantage of this process was to make negotiating the contract easier, reducing public expense, and speeding up construction, while minimising such design work as remained for the contractor. All the member companies of the Eiffage group had some role in the construction work. The construction consortium was made up of the Eiffage TP company for the concrete part, the Eiffel company for the steel roadway (Gustave Eiffel built the Garabit viaduct in 1884, a railway bridge in the neighbouring Cantal département), and the Enerpac company for the roadway's hydraulic supports. The engineering group Setec has authority in the project, with SNCF engineering having partial control. was responsible for the job of the bituminous road surface on the bridge deck, and Forclum (fr) for electrical installations. Management was handled by Eiffage Concessions. The only other business that had a notable role on the building site was Freyssinet, a subsidiary of the Vinci Group specialising in prestressing. It installed the cable stays and put them under tension, while the prestress division of Eiffage was responsible for prestressing the pillar heads. The steel road deck, and the hydraulic action of the road deck were designed by the Walloon engineering firm Greisch from Liège, Belgium, also an information and communication technologies (ICT) company of the Walloon Region. They carried out the general calculations and the resistance calculations for winds of up to . They also applied the launching technology. The sliding shutter technology for the bridge piers came from PERI. Costs and resources The bridge's construction cost up to , with a toll plaza north of the viaduct, costing an additional . The builders, Eiffage, financed the construction in return for a concession to collect the tolls for 75 years, until 2080. However, if the concession yields high revenues, the French government can assume control of the bridge as early as 2044. The project required about of concrete, of steel for the reinforced concrete, and of pre-stressed steel for the cables and shrouds. The builder claims that the lifetime of the bridge will be at least 120 years. Opposition Numerous organisations opposed the project, including the World Wildlife Fund (WWF), France Nature Environnement, the national federation of motorway users, and Environmental Action. Opponents advanced several arguments: The westernmost route would be better, longer by , but a third of the cost with its three more conventional structures. The objective of the viaduct would not be achieved; because of the toll, the viaduct would be little used, and the project would not solve Millau's congestion problems. The project would never break even; toll income would never amortise the initial investment, and the contractor would have to be supported by subsidies. The technical difficulties were too great, and the bridge would be dangerous and unsustainable; the pylons, sitting on the shale of the Tarn Valley, would not support the structure adequately. The viaduct represented a detour, reducing the number of visitors passing through Millau and slowing its economy. Construction Two weeks after the laying of the first stone on 14 December 2001, workers started digging deep shafts for the pilings. Each pylon is supported by four concrete pilings. Each piling is deep and in diameter, assuring the stability of the pylons. At the top of the pilings a large footing was poured, in thickness,to reinforce the strength of the pilings. The of concrete necessary for the footings was poured at the same time as pilings. In March 2002, the pylons emerged from the ground. The speed of construction then rapidly increased. Every three days, each pylon increased in height by . This performance was mainly due to sliding shuttering. Thanks to a system of shoe anchorages and fixed rails in the heart of the pylons, a new layer of concrete could be poured every 20 minutes. Launching The bridge road deck was constructed on plateaus at both ends of the viaduct, and pushed onto the pylons using bridge launching techniques. Each half of the assembled road deck was pushed lengthwise from the plateaus to the pylons, passing across one pylon to the next. During the launching, the road deck was also supported by eight temporary towers, which were removed near the end of construction. In addition to hydraulic jacks on each plateau pushing the road decks, each pylon was topped with a mechanism on top of each pylon that also pushed the deck. This mechanism consisted of a computer-controlled pair of wedges under the deck manipulated by hydraulics. The upper and lower wedge of each pair pointed in opposite directions. The wedges were hydraulically operated, and moved repeatedly in the following sequence: The lower wedge slides under the upper wedge, raising it to the roadway above, and then forcing the upper wedge still higher to lift the roadway Both wedges move forward together, advancing the roadway a short distance The lower wedge retracts from under the upper wedge, lowering the roadway and allowing the upper wedge to drop away from the roadway; the lower wedge then moves back all the way to its starting position. There is now a linear distance between the two wedges equal to the distance forward the roadway has just moved. The upper wedge moves backward, placing it further back along the roadway, adjacent to the front tip of the lower wedge and ready to repeat the cycle and advance the roadway by another increment. The launching advanced the road deck at per cycle which was roughly four minutes long. The mast pieces were driven over the new road deck lying down horizontally. The pieces were joined to form the one complete mast, still lying horizontally. The mast was then tilted upwards, as one piece, at one time in a tricky operation. In this way, each mast was erected on top of the corresponding concrete pylon. The stays connecting the masts and the deck were then installed, and the bridge was tensioned overall, and weight tested. After this, the temporary pylons could be removed. Timeline 16 October 2001: work begins 14 December 2001: laying of the first stone January 2002: laying pier foundations March 2002: start of work on the pier support C8 June 2002: support C8 completed, start of work on piers July 2002: start of work on the foundations of temporary, height adjustable roadway supports August 2002: start of work on pier support C0 September 2002: assembly of roadway begins November 2002: first piers complete 25–26 February 2003: laying of first pieces of roadway November 2003: completion of the last piers (piers P2 at and P3 at are the highest piers in the world) 28 May 2004: the pieces of roadway are several centimetres apart, their juncture to be accomplished within two weeks 2nd half of 2004: installation of the pylons and shrouds, removal of the temporary roadway supports 14 December 2004: official inauguration 16 December 2004: opening of the viaduct, ahead of schedule 10 January 2005: initial planned opening date Construction records The construction Millau Viaduct broke several records: A mega-structure bridge constructed several hundred meters above ground, without any loss of life throughout the three year construction period (First stone laid on 14 December 2001, Official inauguration on 14 December 2004). The highest pylons in the world: pylons P2 and P3, and in height respectively, broke the French record previously held by the Tulle and Verrières viaducts (), and the world record previously held by the Kochertal Viaduct (Germany), which is at its highest; The highest bridge tower in the world: the mast atop pylon P2 peaks at ; The highest road bridge deck in Europe, above the Tarn at its highest point; it is nearly twice as tall as the previous tallest vehicular bridges in Europe, the Europabrücke in Austria and the Italia Viaduct in Italy. Since opening in 2004, the deck height of Millau has been surpassed by several suspension bridges in China, including Sidu River Bridge, Baling River Bridge, and two spans (Beipan River Guanxing Highway Bridge and Beipan River Hukun Expressway Bridge) over the Beipan River. In 2012, Mexico's Baluarte Bridge surpassed Millau as the world's highest cable-stayed bridge. The Royal Gorge suspension bridge in the U.S. state of Colorado is also higher, with a bridge deck approximately over the Arkansas River. Location The Millau Viaduct is on the territory of the communes of Millau and Creissels, France, in the département of Aveyron. Before the bridge was constructed, traffic had to descend into the Tarn valley and pass along the route nationale N9 near the town of Millau, causing much traffic congestion at the beginning and end of the July and August holiday season. The bridge now traverses the Tarn valley above its lowest point, linking two limestone plateaus, the Causse du Larzac and the , and is inside the perimeter of the Grands Causses regional natural park. The Millau Viaduct forms the last link of the existing A75 autoroute (known as "la Méridienne"), from Clermont-Ferrand to Béziers. The A75, with the A10 and A71, provides a continuous high-speed route south from Paris through Clermont-Ferrand to the Languedoc region, thence to Spain, considerably reducing the cost and time of vehicle traffic travelling along this route. Many tourists heading to southern France and Spain follow this route because it is direct and without tolls for the between Clermont-Ferrand and Béziers, except for the bridge. The Eiffage group, which constructed the Viaduct also operates it, under a government contract, which allows the company to collect tolls for up to 75 years. As of 2018, the toll bridge costs for light automobiles (or during the peak season of 15 June to 15 September). Structure Pylons and abutments Each of the seven pylons is supported by four deep shafts, deep and in diameter. The abutments are concrete structures that provide anchorage for the road deck to the ground in the Causse du Larzac and the Causse Rouge. Road deck The metallic road deck, which appears very light despite its total mass of around , is long and wide. It comprises eight spans. The six central spans measure , and the two outer spans are . These are composed of 173 central box beams, the spinal column of the construction, onto which the lateral floors and the lateral box beams were welded. The central box beams have a cross-section, and a length of for a total weight of . The deck has an inverse airfoil shape, providing negative lift in strong wind conditions. Masts The seven masts, each high, and weighing around , are set on top of the concrete pylons. Between each of them, eleven stays (steel cables) are anchored, providing support for the road deck. Cable stays Each mast of the Viaduct is equipped with a monoaxial layer of eleven pairs of cable-stays; laid face to face. Depending on their length, the cable stays were made of 55 to 91 high tensile steel cables, or strands, themselves formed of seven strands of steel (a central strand with six intertwined strands). Each strand has triple protection against corrosion (galvanisation, a coating of petroleum wax, and an extruded polyethylene sheath). The exterior envelope of the stays is itself coated along its entire length with a double helical weatherstrip. The idea is to avoid running water which, in high winds, could cause vibration in the stays and compromise the stability of the viaduct. The stays were installed by the Freyssinet company. Road surface To allow for deformations of the metal road deck under traffic, a special surface of modified bitumen was installed by research teams from . The surface is somewhat flexible to adapt to deformations in the steel deck without cracking, but it must nevertheless have sufficient strength to withstand motorway conditions (fatigue, density, texture, adherence, anti-rutting etc.). The 'ideal formula' was found after two years of research. Electrical installations The electrical installations of the viaduct are large in proportion to the size of the bridge. There are of high-current cables, of fibre optics, of low-current cables, and 357 telephone sockets; allowing maintenance teams to communicate with each other and with the command post. These are situated on the deck, on the pylons, and on the masts. The pylons, road deck, masts, and cable stays are equipped with a multitude of sensors to enable structural health monitoring. These are designed to detect the slightest movement in the Viaduct, and measure its resistance to wear-and-tear over time. Anemometers, accelerometers, inclinometers, and temperature sensors are all used for the instrumentation network. Twelve fibre optic extensometers were installed in the base of pylon P2. Being the tallest of all, it is therefore under the most intense stress. These sensors detect movements on the order of a micrometre. Other extensometers, electrical this time, are distributed on top of P2 and P7. This apparatus is capable of taking up to 100 readings per second. In high winds, they continuously monitor the reactions of the Viaduct to extreme conditions. Accelerometers placed strategically on the road deck monitor the oscillations that can affect the metal structure. Displacements of the deck on the abutment level are measured to the nearest millimetre. The cable stays are also instrumented, and their ageing meticulously analysed. Additionally, two piezoelectric sensors gather traffic data: weight of vehicles, average speed, density of the flow of traffic, etc. This system can distinguish between fourteen different types of vehicle. The data is transmitted by an Ethernet network to a computer in the IT room at the management building situated near the toll plaza. Toll plaza The toll plaza is on the A75 autoroute; the bridge toll booths and the buildings for the commercial and technical management teams are situated north of the viaduct. The toll plaza is protected by a canopy in the shape of a leaf, formed from tendrilled concrete, using the ceracem process. Consisting of 53 elements (voussoirs), the canopy is long and wide. It weighs around . The toll plaza can accommodate sixteen lanes of traffic, eight in each direction. At times of low traffic volume, the central booth is capable of servicing vehicles in both directions. A car park and viewing station, equipped with public toilets, is situated at each side of the toll plaza. The total cost was . Rest area of Brocuéjouls The rest area of Brocuéjouls, named Aire du Viaduc de Millau, is situated just north of the viaduct, and is centred on an old farm building named 'Ferme de Brocuéjouls'. It was inaugurated by the prefect of Aveyron, Chantal Jourdan, on 30 June 2006, after 7 months of works. The farm and its surroundings can accommodate entertainment and tourism promotion activities. The cost of this work amounted to : of state funds for the realisation of the area (access roads, parking, rest area, toilets, etc.) for the restoration of the old farm building of Brocuéjouls (all two tranches) Statistics : total length of the roadway 7: number of piers : height of Pier 7, the shortest : height of Pier 2, the tallest ( at the roadway's level) : height of a mast 154: number of shrouds : average height of the roadway : thickness of the roadway : width of the roadway : total volume of concrete used : total weight of the bridge 10,000–25,000 vehicles: estimated daily traffic : typical automobile toll (price increasing in summer), : horizontal radius of curvature of the road deck Impact and events Pedestrian sporting events Unusually for a bridge closed to pedestrians, a run took place in 2004, and another on 13 May 2007: December 2004 – 19,000 walkers and runners of the Three Bridge Walk had the privilege of crossing the bridge deck for the first time, but the walk was not authorised to go further than pylon P1; the bridge was still closed to traffic. 13 May 2007 – 10,496 runners took the departure of the race from Place de Mandarous, in the centre of Millau, to the southern end of the Viaduct. After starting on the northern side, they crossed the viaduct, then retraced their steps. Total distance: . Events and popular culture In 2004, a fire started on the slope of the because of a spark originating from a welder; some trees were destroyed in the fire. The speed limit on the bridge was reduced from to because tourists were slowing down to take photos. Soon after the bridge opened to traffic, cars were stopping on the hard shoulder so that travelers could view the landscape and the bridge. A postage stamp was designed by Sarah Lazarevic to commemorate the opening of the Viaduct. The Chinese transport minister at the time visited the bridge on the first anniversary of its opening. The commission was impressed by the technical prowess of the bridge's immense construction, but also by the legal and financial assembly of the Viaduct. However, according to the minister, he did not envisage building a counterpart in People's Republic of China. The cabinet of the governor of California Arnold Schwarzenegger, who envisaged the construction of a bridge in San Francisco Bay, asked the council of the town hall of Millau about the popularity of the construction of the viaduct. This bridge was featured in a scene of the 2007 film Mr. Bean's Holiday. The bridge also appeared in the ending of the console and PC versions of the third-person shooter video game James Bond 007: Blood Stone, where James Bond confronts the game's true antagonist. The hosts of the British motoring show Top Gear featured the bridge during Series 7, when they took a Ford GT, Pagani Zonda, and Ferrari F430 Spyder on a road trip across France to see the newly completed bridge. Richard Hammond, one of the above hosts on Top Gear, explored the engineering aspects in the construction of the Millau Viaduct in Series 2 of Richard Hammond's Engineering Connections. The bridge was featured in Series 2 of World's Greatest Bridges. Construction of the bridge was featured in the series How Did They Build That?	Technology	Bridges	null
681218	https://en.wikipedia.org/wiki/Omega-6%20fatty%20acid	Omega-6 fatty acid	Omega−6 fatty acids (also referred to as ω−6 fatty acids or n−6 fatty acids) are a family of polyunsaturated fatty acids (PUFA) that share a final carbon-carbon double bond in the n−6 position, that is, the sixth bond, counting from the methyl end. Health effects The American Heart Association "supports an omega-6 PUFA intake of at least 5% to 10% of energy in the context of other AHA lifestyle and dietary recommendations. To reduce omega-6 PUFA intakes from their current levels would be more likely to increase than to decrease risk for coronary heart disease." A 2018 review found that an increased intake of omega−6 fatty acids reduces total serum cholesterol and may reduce myocardial infarction (heart attack), but found no significant change in LDL cholesterol and triglycerides. A 2021 review found that omega−6 supplements do not affect the risk of CVD morbidity and mortality. A 2023 review found that omega−6 polyunsaturated fatty acids are associated with lower risk of high blood pressure. Omega−6 fatty acids are not associated with atrial fibrillation. A review and meta-analysis of observational studies by the World Health Organization (WHO) found that higher intakes of omega-6 are associated with a 9% reduced risk of all-cause mortality and a 31% increased risk of postmenopausal breast cancer. The increased risk of breast cancer has not been confirmed in randomized controlled trials. A scoping review for Nordic Nutrition Recommendations 2023 found that partial replacement of saturated fatty acid with omega-6 fatty acid decreases risk of cardiovascular disease and improves the blood lipid profile. Dietary sources Dietary sources of omega−6 fatty acids include: poultry eggs nuts hulled sesame seeds cereals durum wheat whole-grain breads pumpkin seeds hemp seeds Vegetable oils Vegetable oils are a major source of omega−6 linoleic acid. Worldwide, more than 100 million metric tons of vegetable oils are extracted annually from palm fruits, soybean seeds, grape seeds, and sunflower seeds, providing more than 32 million metric tons of omega−6 linoleic acid and 4 million metric tons of omega−3 alpha-linolenic acid. List of omega−6 fatty acids The melting point of the fatty acids increases as the number of carbons in the chain increases.	Biology and health sciences	Lipids	Biology
681579	https://en.wikipedia.org/wiki/Quasiparticle	Quasiparticle	In condensed matter physics, a quasiparticle is a concept used to describe a collective behavior of a group of particles that can be treated as if they were a single particle. Formally, quasiparticles and collective excitations are closely related phenomena that arise when a microscopically complicated system such as a solid behaves as if it contained different weakly interacting particles in vacuum. For example, as an electron travels through a semiconductor, its motion is disturbed in a complex way by its interactions with other electrons and with atomic nuclei. The electron behaves as though it has a different effective mass travelling unperturbed in vacuum. Such an electron is called an electron quasiparticle. In another example, the aggregate motion of electrons in the valence band of a semiconductor or a hole band in a metal behave as though the material instead contained positively charged quasiparticles called electron holes. Other quasiparticles or collective excitations include the phonon, a quasiparticle derived from the vibrations of atoms in a solid, and the plasmons, a particle derived from plasma oscillation. These phenomena are typically called quasiparticles if they are related to fermions, and called collective excitations if they are related to bosons, although the precise distinction is not universally agreed upon. Thus, electrons and electron holes (fermions) are typically called quasiparticles, while phonons and plasmons (bosons) are typically called collective excitations. The quasiparticle concept is important in condensed matter physics because it can simplify the many-body problem in quantum mechanics. The theory of quasiparticles was started by the Soviet physicist Lev Landau in the 1930s. Overview General introduction Solids are made of only three kinds of particles: electrons, protons, and neutrons. None of these are quasiparticles; instead a quasiparticle is an emergent phenomenon that occurs inside the solid. Therefore, while it is quite possible to have a single particle (electron, proton, or neutron) floating in space, a quasiparticle can only exist inside interacting many-particle systems such as solids. Motion in a solid is extremely complicated: Each electron and proton is pushed and pulled (by Coulomb's law) by all the other electrons and protons in the solid (which may themselves be in motion). It is these strong interactions that make it very difficult to predict and understand the behavior of solids (see many-body problem). On the other hand, the motion of a non-interacting classical particle is relatively simple; it would move in a straight line at constant velocity. This is the motivation for the concept of quasiparticles: The complicated motion of the real particles in a solid can be mathematically transformed into the much simpler motion of imagined quasiparticles, which behave more like non-interacting particles. In summary, quasiparticles are a mathematical tool for simplifying the description of solids. Relation to many-body quantum mechanics The principal motivation for quasiparticles is that it is almost impossible to directly describe every particle in a macroscopic system. For example, a barely-visible (0.1mm) grain of sand contains around 1017 nuclei and 1018 electrons. Each of these attracts or repels every other by Coulomb's law. In principle, the Schrödinger equation predicts exactly how this system will behave. But the Schrödinger equation in this case is a partial differential equation (PDE) on a 3×1018-dimensional vector space—one dimension for each coordinate (x, y, z) of each particle. Directly and straightforwardly trying to solve such a PDE is impossible in practice. Solving a PDE on a 2-dimensional space is typically much harder than solving a PDE on a 1-dimensional space (whether analytically or numerically); solving a PDE on a 3-dimensional space is significantly harder still; and thus solving a PDE on a 3×1018-dimensional space is quite impossible by straightforward methods. One simplifying factor is that the system as a whole, like any quantum system, has a ground state and various excited states with higher and higher energy above the ground state. In many contexts, only the "low-lying" excited states, with energy reasonably close to the ground state, are relevant. This occurs because of the Boltzmann distribution, which implies that very-high-energy thermal fluctuations are unlikely to occur at any given temperature. Quasiparticles and collective excitations are a type of low-lying excited state. For example, a crystal at absolute zero is in the ground state, but if one phonon is added to the crystal (in other words, if the crystal is made to vibrate slightly at a particular frequency) then the crystal is now in a low-lying excited state. The single phonon is called an elementary excitation. More generally, low-lying excited states may contain any number of elementary excitations (for example, many phonons, along with other quasiparticles and collective excitations). When the material is characterized as having "several elementary excitations", this statement presupposes that the different excitations can be combined. In other words, it presupposes that the excitations can coexist simultaneously and independently. This is never exactly true. For example, a solid with two identical phonons does not have exactly twice the excitation energy of a solid with just one phonon, because the crystal vibration is slightly anharmonic. However, in many materials, the elementary excitations are very close to being independent. Therefore, as a starting point, they are treated as free, independent entities, and then corrections are included via interactions between the elementary excitations, such as "phonon-phonon scattering". Therefore, using quasiparticles / collective excitations, instead of analyzing 1018 particles, one needs to deal with only a handful of somewhat-independent elementary excitations. It is, therefore, an effective approach to simplify the many-body problem in quantum mechanics. This approach is not useful for all systems, however. For example, in strongly correlated materials, the elementary excitations are so far from being independent that it is not even useful as a starting point to treat them as independent. Distinction between quasiparticles and collective excitations Usually, an elementary excitation is called a "quasiparticle" if it is a fermion and a "collective excitation" if it is a boson. However, the precise distinction is not universally agreed upon. There is a difference in the way that quasiparticles and collective excitations are intuitively envisioned. A quasiparticle is usually thought of as being like a dressed particle: it is built around a real particle at its "core", but the behavior of the particle is affected by the environment. A standard example is the "electron quasiparticle": an electron in a crystal behaves as if it had an effective mass which differs from its real mass. On the other hand, a collective excitation is usually imagined to be a reflection of the aggregate behavior of the system, with no single real particle at its "core". A standard example is the phonon, which characterizes the vibrational motion of every atom in the crystal. However, these two visualizations leave some ambiguity. For example, a magnon in a ferromagnet can be considered in one of two perfectly equivalent ways: (a) as a mobile defect (a misdirected spin) in a perfect alignment of magnetic moments or (b) as a quantum of a collective spin wave that involves the precession of many spins. In the first case, the magnon is envisioned as a quasiparticle, in the second case, as a collective excitation. However, both (a) and (b) are equivalent and correct descriptions. As this example shows, the intuitive distinction between a quasiparticle and a collective excitation is not particularly important or fundamental. The problems arising from the collective nature of quasiparticles have also been discussed within the philosophy of science, notably in relation to the identity conditions of quasiparticles and whether they should be considered "real" by the standards of, for example, entity realism. Effect on bulk properties By investigating the properties of individual quasiparticles, it is possible to obtain a great deal of information about low-energy systems, including the flow properties and heat capacity. In the heat capacity example, a crystal can store energy by forming phonons, and/or forming excitons, and/or forming plasmons, etc. Each of these is a separate contribution to the overall heat capacity. History The idea of quasiparticles originated in Lev Landau's theory of Fermi liquids, which was originally invented for studying liquid helium-3. For these systems a strong similarity exists between the notion of quasiparticle and dressed particles in quantum field theory. The dynamics of Landau's theory is defined by a kinetic equation of the mean-field type. A similar equation, the Vlasov equation, is valid for a plasma in the so-called plasma approximation. In the plasma approximation, charged particles are considered to be moving in the electromagnetic field collectively generated by all other particles, and hard collisions between the charged particles are neglected. When a kinetic equation of the mean-field type is a valid first-order description of a system, second-order corrections determine the entropy production, and generally take the form of a Boltzmann-type collision term, in which figure only "far collisions" between virtual particles. In other words, every type of mean-field kinetic equation, and in fact every mean-field theory, involves a quasiparticle concept. Common examples This section contains most commmon examples of quasiparticles and collective excitations. In solids, an electron quasiparticle is an electron as affected by the other forces and interactions in the solid. The electron quasiparticle has the same charge and spin as a "normal" (elementary particle) electron, and like a normal electron, it is a fermion. However, its mass can differ substantially from that of a normal electron; see the article effective mass. Its electric field is also modified, as a result of electric field screening. In many other respects, especially in metals under ordinary conditions, these so-called Landau quasiparticles closely resemble familiar electrons; as Crommie's "quantum corral" showed, an STM can image their interference upon scattering. A hole is a quasiparticle consisting of the lack of an electron in a state; it is most commonly used in the context of empty states in the valence band of a semiconductor. A hole has the opposite charge of an electron. A phonon is a collective excitation associated with the vibration of atoms in a rigid crystal structure. It is a quantum of a sound wave. A magnon is a collective excitation associated with the electrons' spin structure in a crystal lattice. It is a quantum of a spin wave. In materials, a photon quasiparticle is a photon as affected by its interactions with the material. In particular, the photon quasiparticle has a modified relation between wavelength and energy (dispersion relation), as described by the material's index of refraction. It may also be termed a polariton, especially near a resonance of the material. For example, an exciton-polariton is a superposition of an exciton and a photon; a phonon-polariton is a superposition of a phonon and a photon. A plasmon is a collective excitation, which is the quantum of plasma oscillations (wherein all the electrons simultaneously oscillate with respect to all the ions). A polaron is a quasiparticle which comes about when an electron interacts with the polarization of its surrounding ions. An exciton is an electron and hole bound together.	Physical sciences	States of matter	Physics
681666	https://en.wikipedia.org/wiki/Penrose%20diagram	Penrose diagram	In theoretical physics, a Penrose diagram (named after mathematical physicist Roger Penrose) is a two-dimensional diagram capturing the causal relations between different points in spacetime through a conformal treatment of infinity. It is an extension (suitable for the curved spacetimes of e.g. general relativity) of the Minkowski diagram of special relativity where the vertical dimension represents time, and the horizontal dimension represents a space dimension. Using this design, all light rays take a 45° path . Locally, the metric on a Penrose diagram is conformally equivalent to the metric of the spacetime depicted. The conformal factor is chosen such that the entire infinite spacetime is transformed into a Penrose diagram of finite size, with infinity on the boundary of the diagram. For spherically symmetric spacetimes, every point in the Penrose diagram corresponds to a 2-dimensional sphere . Basic properties While Penrose diagrams share the same basic coordinate vector system of other spacetime diagrams for local asymptotically flat spacetime, it introduces a system of representing distant spacetime by shrinking or "triturando" distances that are further away. Straight lines of constant time and straight lines of constant space coordinates therefore become hyperbolae, which appear to converge at points in the corners of the diagram. These points and boundaries represent conformal infinity for spacetime, which was first introduced by Penrose in 1963. Penrose diagrams are more properly (but less frequently) called Penrose–Carter diagrams (or Carter–Penrose diagrams), acknowledging both Brandon Carter and Roger Penrose, who were the first researchers to employ them. They are also called conformal diagrams, or simply spacetime diagrams (although the latter may refer to Minkowski diagrams). Two lines drawn at 45° angles should intersect in the diagram only if the corresponding two light rays intersect in the actual spacetime. So, a Penrose diagram can be used as a concise illustration of spacetime regions that are accessible to observation. The diagonal boundary lines of a Penrose diagram correspond to the region called "null infinity", or to singularities where light rays must end. Thus, Penrose diagrams are also useful in the study of asymptotic properties of spacetimes and singularities. An infinite static Minkowski universe, coordinates is related to Penrose coordinates by: The corners of the Penrose diagram, which represent the spacelike and timelike conformal infinities, are from the origin. Black holes Penrose diagrams are frequently used to illustrate the causal structure of spacetimes containing black holes. Singularities in the Schwarzschild solution are denoted by a spacelike boundary, unlike the timelike boundary found on conventional spacetime diagrams. This is due to the interchanging of timelike and spacelike coordinates within the horizon of a black hole (since space is uni-directional within the horizon, just as time is uni-directional outside the horizon). The singularity is represented by a spacelike boundary to make it clear that once an object has passed the horizon it will inevitably hit the singularity even if it attempts to take evasive action. Penrose diagrams are often used to illustrate the hypothetical Einstein–Rosen bridge connecting two separate universes in the maximally extended Schwarzschild black hole solution. The precursors to the Penrose diagrams were Kruskal–Szekeres diagrams. (The Penrose diagram adds to Kruskal and Szekeres' diagram the conformal crunching of the regions of flat spacetime far from the hole.) These introduced the method of aligning the event horizon into past and future horizons oriented at 45° angles (since one would need to travel faster than light to cross from the Schwarzschild radius back into flat spacetime); and splitting the singularity into past and future horizontally-oriented lines (since the singularity "cuts off" all paths into the future once one enters the hole). The Einstein–Rosen bridge closes off (forming "future" singularities) so rapidly that passage between the two asymptotically flat exterior regions would require faster-than-light velocity, and is therefore impossible. In addition, highly blue-shifted light rays (called a blue sheet) would make it impossible for anyone to pass through. The maximally extended solution does not describe a typical black hole created from the collapse of a star, as the surface of the collapsed star replaces the sector of the solution containing the past-oriented white hole geometry and other universe. While the basic space-like passage of a static black hole cannot be traversed, the Penrose diagrams for solutions representing rotating and/or electrically charged black holes illustrate these solutions' inner event horizons (lying in the future) and vertically oriented singularities, which open up what is known as a time-like "wormhole" allowing passage into future universes. In the case of the rotating hole, there is also a "negative" universe entered through a ring-shaped singularity (still portrayed as a line in the diagram) that can be passed through if entering the hole close to its axis of rotation. These features of the solutions are, however, not stable under perturbations and not believed to be a realistic description of the interior regions of such black holes; the true character of their interiors is still an open question.	Physical sciences	Theory of relativity	Physics
681875	https://en.wikipedia.org/wiki/General%20covariance	General covariance	In theoretical physics, general covariance, also known as diffeomorphism covariance or general invariance, consists of the invariance of the form of physical laws under arbitrary differentiable coordinate transformations. The essential idea is that coordinates do not exist a priori in nature, but are only artifices used in describing nature, and hence should play no role in the formulation of fundamental physical laws. While this concept is exhibited by general relativity, which describes the dynamics of spacetime, one should not expect it to hold in less fundamental theories. For matter fields taken to exist independently of the background, it is almost never the case that their equations of motion will take the same form in curved space that they do in flat space. Overview A physical law expressed in a generally covariant fashion takes the same mathematical form in all coordinate systems, and is usually expressed in terms of tensor fields. The classical (non-quantum) theory of electrodynamics is one theory that has such a formulation. Albert Einstein proposed this principle for his special theory of relativity; however, that theory was limited to spacetime coordinate systems related to each other by uniform inertial motion, meaning relative motion in any straight line without acceleration. Einstein recognized that the general principle of relativity should also apply to accelerated relative motions, and he used the newly developed tool of tensor calculus to extend the special theory's global Lorentz covariance (applying only to inertial frames) to the more general local Lorentz covariance (which applies to all frames), eventually producing his general theory of relativity. The local reduction of the metric tensor to the Minkowski metric tensor corresponds to free-falling (geodesic) motion, in this theory, thus encompassing the phenomenon of gravitation. Much of the work on classical unified field theories consisted of attempts to further extend the general theory of relativity to interpret additional physical phenomena, particularly electromagnetism, within the framework of general covariance, and more specifically as purely geometric objects in the spacetime continuum. Remarks The relationship between general covariance and general relativity may be summarized by quoting a standard textbook: A more modern interpretation of the physical content of the original principle of general covariance is that the Lie group GL4(R) is a fundamental "external" symmetry of the world. Other symmetries, including "internal" symmetries based on compact groups, now play a major role in fundamental physical theories.	Physical sciences	Theory of relativity	Physics
682382	https://en.wikipedia.org/wiki/Phage%20therapy	Phage therapy	Phage therapy, viral phage therapy, or phagotherapy is the therapeutic use of bacteriophages for the treatment of pathogenic bacterial infections. This therapeutic approach emerged at the beginning of the 20th century but was progressively replaced by the use of antibiotics in most parts of the world after the Second World War. Bacteriophages, known as phages, are a form of virus that attach to bacterial cells and inject their genome into the cell. The bacteria's production of the viral genome interferes with its ability to function, halting the bacterial infection. The bacterial cell causing the infection is unable to reproduce and instead produces additional phages. Phages are very selective in the strains of bacteria they are effective against. Advantages include reduced side effects and reduced risk of the bacterium developing resistance, since bacteriophages are much more specific than antibiotics. They are typically harmless not only to the host organism but also to other beneficial bacteria, such as the gut microbiota, reducing the chances of opportunistic infections. They have a high therapeutic index; that is, phage therapy would be expected to give rise to few side effects, even at higher-than-therapeutic levels. Because phages replicate in vivo (in cells of living organism), a smaller effective dose can be used. Disadvantages include the difficulty of finding an effective phage for a particular infection; a phage will kill a bacterium only if it matches the specific strain. However, virulent phages can be isolated much more easily than other compounds and natural products. Consequently, phage mixtures ("cocktails") are sometimes used to improve the chances of success. Alternatively, samples taken from recovering patients sometimes contain appropriate phages that can be grown to cure other patients infected with the same strain. Ongoing challenges include the need to increase phage collections from reference phage banks, the development of efficient phage screening methods for the fast identification of the therapeutic phage(s), the establishment of efficient phage therapy strategies to tackle infectious biofilms, the validation of feasible phage production protocols that assure quality and safety of phage preparations, and the guarantee of stability of phage preparations during manufacturing, storage, and transport. Phages tend to be more successful than antibiotics where there is a biofilm covered by a polysaccharide layer, which antibiotics typically cannot penetrate. Phage therapy can disperse the biofilm generated by antibiotic-resistant bacteria. However, the interactions between phages and biofilms can be complex, with phages developing symbiotic as well as predatory relationships with biofilms. Phages are currently being used therapeutically to treat bacterial infections that do not respond to conventional antibiotics, particularly in Russia and Georgia. There is also a phage therapy unit in Wrocław, Poland, established in 2005, which continues several-decades-long research by the Institute of Immunology and Experimental Therapy of the Polish Academy of Sciences, the only such centre in a European Union country. Phages are the subject of renewed clinical attention in Western countries, such as the United States. In 2019, the United States Food and Drug Administration approved the first US clinical trial for intravenous phage therapy. Phage therapy has many potential applications in human medicine as well as dentistry, veterinary science, and agriculture. If the target host of a phage therapy treatment is not an animal, the term "biocontrol" (as in phage-mediated biocontrol of bacteria) is usually employed, rather than "phage therapy". History The discovery of bacteriophages was reported by British bacteriologist Frederick Twort in 1915 and by French microbiologist Felix d'Hérelle in 1917. D'Hérelle said that the phages always appeared in the stools of Shigella dysentery patients shortly before they began to recover. He "quickly learned that bacteriophages are found wherever bacteria thrive: in sewers, in rivers that catch waste runoff from pipes, and in the stools of convalescent patients". Phage therapy was immediately recognized by many to be a key way forward for the eradication of pathogenic bacterial infections. A Georgian, George Eliava, was making similar discoveries. He travelled to the Pasteur Institute in Paris, where he met d'Hérelle, and in 1923, he founded the Institute of Bacteriology, which later became known as the George Eliava Institute, in Tbilisi, Georgia, devoted to the development of phage therapy. Phage therapy is used in Russia, Georgia and Poland, and was used prophylactically for a time in the Soviet army, most notably during the Second World War. In Russia, extensive research and development soon began in this field. In the United States during the 1940s, commercialization of phage therapy was undertaken by Eli Lilly and Company. While knowledge was being accumulated regarding the biology of phages and how to use phage cocktails correctly, early uses of phage therapy were often unreliable. Since the early 20th century, research into the development of viable therapeutic antibiotics had also been underway, and by 1942, the antibiotic penicillin G had been successfully purified and saw use during the Second World War. The drug proved to be extraordinarily effective in the treatment of injured Allied soldiers whose wounds had become infected. By 1944, large-scale production of penicillin had been made possible, and in 1945, it became publicly available in pharmacies. Due to the drug's success, it was marketed widely in the US and Europe, leading Western scientists to mostly lose interest in further use and study of phage therapy for some time. Isolated from Western advances in antibiotic production in the 1940s, Russian scientists continued to develop already successful phage therapy to treat the wounds of soldiers in field hospitals. During World War II, the Soviet Union used bacteriophages to treat soldiers infected with various bacterial diseases, such as dysentery and gangrene. Russian researchers continued to develop and to refine their treatments and to publish their research and results. However, due to the scientific barriers of the Cold War, this knowledge was not translated and did not proliferate across the world. A summary of these publications was published in English in 2009 in "A Literature Review of the Practical Application of Bacteriophage Research". There is an extensive library and research center at the George Eliava Institute in Tbilisi, Georgia. Phage therapy is today a widespread form of treatment in that region. As a result of the development of antibiotic resistance since the 1950s and an advancement of scientific knowledge, there has been renewed interest worldwide in the ability of phage therapy to eradicate bacterial infections and chronic polymicrobial biofilm (including in industrial situations). Phages have been investigated as a potential means to eliminate pathogens like Campylobacter in raw food and Listeria in fresh food or to reduce food spoilage bacteria. In agricultural practice, phages have been used to fight pathogens like Campylobacter, Escherichia, and Salmonella in farm animals, Lactococcus and Vibrio pathogens in fish aquaculture, and Erwinia, Xanthomonas, and others in plants of agricultural importance. The oldest use is, however, in human medicine. Phages have been used against diarrheal diseases caused by E. coli, Shigella, or Vibrio and against wound infections caused by facultative pathogens of the skin like staphylococci and streptococci. Recently, the phage therapy approach has been applied to systemic and even intracellular infections, and non-replicating phage and isolated phage enzymes like lysins have been added to the antimicrobial arsenal. However, actual proof for the efficacy of these phage approaches in the field or the hospital is not available. Some of the interest in the West can be traced back to 1994, when James Soothill demonstrated (in an animal model) that the use of phages could improve the success of skin grafts by reducing the underlying Pseudomonas aeruginosa infection. Recent studies have provided additional support for these findings in the model system. Although not "phage therapy" in the original sense, the use of phages as delivery mechanisms for traditional antibiotics constitutes another possible therapeutic use. The use of phages to deliver antitumor agents has also been described in preliminary in vitro experiments for cells in tissue culture. In June 2015, the European Medicines Agency hosted a one-day workshop on the therapeutic use of bacteriophages, and in July 2015, the US National Institutes of Health hosted a two-day workshop titled "Bacteriophage Therapy: An Alternative Strategy to Combat Drug Resistance". In January 2016, phages were used successfully at Yale University by Benjamin Chan to treat a chronic Pseudomonas aeruginosa infection in ophthalmologist Ali Asghar Khodadoust. This successful treatment of a life-threatening infection sparked a resurgence of interest in phage therapy in the United States. In 2017, a pair of genetically engineered phages along with one naturally occurring (so-called "phage Muddy") each from among those catalogued by SEA-PHAGES (Science Education Alliance-Phage Hunters Advancing Genomics and Evolutionary Science) at the Howard Hughes Medical Institute by Graham Hatfull and colleagues, was used by microbiologist James Soothill at Great Ormond Street Hospital for Children in London to treat an antibiotic-resistant bacterial (Mycobacterium abscessus) infection in a young woman with cystic fibrosis. In 2022, two mycobacteriophages were administered intravenously twice daily to a young man with treatment-refractory Mycobacterium abscessus pulmonary infection and severe cystic fibrosis lung disease. Airway cultures for M. abscessus became negative after approximately 100 days of combined phage and antibiotic treatment, and a variety of biomarkers confirmed the therapeutic response. The individual received a bilateral lung transplant after 379 days of treatment, and cultures from the explanted lung tissue confirmed eradication of the bacteria. In a second case, successful treatment of disseminated cutaneous Mycobacterium chelonae was reported with a single phage administered intravenously twice daily in conjunction with antibiotic and surgical management. Potential benefits Bacteriophage treatment offers a possible alternative to conventional antibiotic treatments for bacterial infection. It is conceivable that, although bacteria can develop resistance to phages, the resistance might be easier to overcome than resistance to antibiotics. Viruses, just like bacteria, can evolve resistance to different treatments. Bacteriophages are very specific, targeting only one or a few strains of bacteria. Traditional antibiotics have a more wide-ranging effect, killing both harmful and useful bacteria, such as those facilitating food digestion. The species and strain specificity of bacteriophages makes it unlikely that harmless or useful bacteria will be killed when fighting an infection. A few research groups in the West are engineering a broader-spectrum phage and also a variety of forms of MRSA treatments, including impregnated wound dressings, preventative treatment for burn victims, and phage-impregnated sutures. Enzybiotics are a new development at Rockefeller University that create enzymes from phages. Purified recombinant phage enzymes can be used as separate antibacterial agents in their own right. Phage therapy also has the potential to prevent or treat infectious diseases of corals. This could mitigate the global coral decline. Applications Collection Phages for therapeutic use can be collected from environmental sources that likely contain high quantities of bacteria and bacteriophages, such as effluent outlets, sewage, or even soil. The samples are taken and applied to bacterial cultures that are to be targeted. If the bacteria die, the phages can be grown in liquid cultures. Modes of treatment Phages are "bacterium-specific", and therefore, it is necessary in many cases to take a swab from the patient and culture it prior to treatment. Occasionally, isolation of therapeutic phages can require a few months to complete, but clinics generally keep supplies of phage cocktails for the most common bacterial strains in a geographical area. Phage cocktails are commonly sold in pharmacies in Eastern European countries, such as Russia and Georgia. The composition of bacteriophagic cocktails has been periodically modified to add phages effective against emerging pathogenic strains. Phages in practice are applied orally, topically on infected wounds or spread onto surfaces, or during surgical procedures. Injection is rarely used, avoiding any risks of trace chemical contaminants that may be present from the bacteria amplification stage, and recognizing that the immune system naturally fights against viruses introduced into the bloodstream or lymphatic system. Reviews of phage therapy indicate that more clinical and microbiological research is needed to meet current standards. Clinical trials Funding for phage therapy research and clinical trials is generally insufficient and difficult to obtain, since it is a lengthy and complex process to patent bacteriophage products. Due to the specificity of phages, phage therapy would be most effective as a cocktail injection, a modality generally rejected by the US Food and Drug Administration (FDA). Therefore, researchers and observers have predicted that if phage therapy is to gain traction, the FDA must change its regulatory stance on combination drug cocktails. Public awareness and education about phage therapy are generally limited to scientific or independent research rather than mainstream media. In 2007, phase-1 and 2 clinical trials were completed at the Royal National Throat, Nose and Ear Hospital, London, for Pseudomonas aeruginosa infections (otitis). Phase-1 clinical trials were conducted at the Southwest Regional Wound Care Center of Lubbock, Texas, for a cocktail of phages against P. aeruginosa, Staphylococcus aureus, and Escherichia coli, developed by Intralytix. , a phase-1 and 2 trial of phage therapy against P. aeruginosa wound infection in France and Belgium in 2015–17, was terminated early due to lack of effectiveness. Locus Biosciences has created a cocktail of three CRISPR-modified phages. A 2019 study examined its effectiveness against E. coli in the urinary tract, and a phase-1 trial was completed shortly before March 2021. In February 2019, the FDA approved the first clinical trial of intravenously administered phage therapy in the United States. In July 2020, the FDA approved the first clinical trial of nebulized phage therapy in the United States. This double-blind, placebo-controlled study at Yale University will be focused on treating P. aeruginosa infections in patients with cystic fibrosis. In February 2020, the FDA approved a clinical trial to evaluate bacteriophage therapy in patients with urinary tract infections. The study started in December 2020 and aims to identify ideal bacteriophage treatment regimens based on improvements in disease control rates. In February 2021, the FDA approved a clinical trial to evaluate bacteriophage therapy in patients with chronic prosthetic joint infections (PJI). The study was to begin in October 2022 and be conducted by Adaptive Phage Therapeutics, in collaboration with the Mayo Clinic. Administration As pills If administered as pills, phages can be freeze-dried; this procedure does not reduce efficiency. Temperature stability up to 55 °C and shelf lives of 14 months have been shown for some types of phages in pill form. Liquid Application in liquid form is possible, stored preferably in refrigerated vials. Oral administration works better when an antacid is included, as this increases the number of phages surviving passage through the stomach. Topical administration often involves application to gauzes that are laid on the area to be treated. Liquid bacteriophages are also utilized for local applications, such as wound dressings and topical treatments, as well as external administration, including sprays and rinses. Via nebulizer The July 2020 application for FDA approval for the first clinical trial of nebulized phage therapy in the United States does not specify a particular type of nebulizer, such as a compressor or ultrasound type. Bacteriophages are studied as potential candidates for treating bacterial lung infections, especially those caused by multidrug-resistant (MDR) bacteria. In these studies, bacteriophage solutions are administered via nebulizers, mostly using the compressor type. The stability and viability of phages during nebulization are crucial for their therapeutic efficacy. Current studies focus on whether phages can remain viable and effective when delivered via nebulizers. The choice of nebulizer can impact the stability and delivery efficiency of phages. Compressor nebulizers are commonly used because they generate a fine mist that can reach the lower respiratory tract. In contrast to the compressor nebulizers, the ultrasound nebulizers can impact the viability of bacteriophages. The ultrasonic waves used to generate the aerosol can cause physical damage to the phages, potentially reducing their effectiveness. Preliminary research suggests the high-frequency vibrations and heat generated during the nebulization process can lead to a significant loss of phage activity. Consequently, one of the main challenges is ensuring that the phages remain undamaged during the nebulization process. Studies have shown that phages can be sensitive to the shear forces generated during nebulization. Still, with proper formulation and device selection, it is possible to maintain their viability, as the current research suggests. Successful treatments Phages were used successfully at Yale University by Benjamin Chan to treat a Pseudomonas infection in 2016. Intravenous phage drip therapy was successfully used to treat a patient with multidrug-resistant Acinetobacter baumannii in Thornton Hospital at UC San Diego in 2017. Nebulized phage therapy has been used successfully to treat numerous patients with cystic fibrosis and multidrug-resistant bacteria at Yale University as part of their compassionate use program. In 2019, a Brownsville, Minnesota resident with a longstanding bacterial infection in his knee received a phage treatment at the Mayo Clinic that eliminated the need for amputation of his lower leg. Individualised phage therapy was also successfully used by Robert T. Schooley and others to treat a case of multi-drug-resistant Acinetobacter baumannii in 2015. In 2022, an individually adjusted phage-antibiotic combination as an antimicrobial resistance treatment was demonstrated and described in detail. The scientists called for scaling up the research and for further development of this approach. Treatment of biofilm infections Phage therapy is being used to great effect in the treatment of biofilm infections, especially Pseudomonas aeruginosa and Staphylococcus aureus. From 78 recent cases of treatment of biofilm infections, 96% of patients saw clinical improvement using phage therapy, and 52% of patients saw complete symptom relief or a full expungement of the affecting bacteria. Biofilm infections are very challenging to treat with antibiotics. The biofilm matrix and surrounding bacterial membranes can bind to the antibiotics, preventing them from penetrating the biofilm. The matrix may contain enzymes that deactivate antibiotics. Biofilms also have low metabolic activity, which means antibiotics that target growing processes have much lower efficacy. These factors make phage therapy an enticing option for the treatment of such infections, and there are currently two ways to go about such treatment. The first is to isolate the initial bacteria and make a specific treatment phage to target it, while the second way is to use a combination of more general phages. The advantage of the second method is that it can easily be made commercially available for treatment, although there are some concerns that it may be substantially less effective. Limitations The high bacterial strain specificity of phage therapy may make it necessary for clinics to make different cocktails for treatment of the same infection or disease, because the bacterial components of such diseases may differ from region to region or even person to person. In addition, this means that "banks" containing many different phages must be kept and regularly updated with new phages. Further, bacteria can evolve different receptors either before or during treatment. This can prevent phages from completely eradicating them. The need for banks of phages makes regulatory testing for safety harder and more expensive under current rules in most countries. Such a process would make the large-scale use of phage therapy difficult. Additionally, patent issues (specifically on living organisms) may complicate distribution for pharmaceutical companies wishing to have exclusive rights over their "invention", which would discourage a commercial corporation from investing capital in this. As has been known for at least thirty years, mycobacteria such as Mycobacterium tuberculosis have specific bacteriophages. No lytic phage has yet been discovered for Clostridioides difficile, which is responsible for many nosocomial diseases, but some temperate phages (integrated in the genome, also called lysogenic) are known for this species; this opens encouraging avenues but with additional risks, as discussed below. The negative public perception of viruses may contribute to the reluctance to embrace phage therapy. Development of resistance One of the major concerns usually associated with phage therapy is the emergence of bacteriophage-insensitive mutants (BIMs) that could hinder the success of this therapy. Several in vitro studies have reported a fast emergence of BIMs within a short time after phage treatment. The emergence of BIMs has also been observed in vivo using different animal models, although this usually occurs later than in vitro (reviewed in ). This fast adaptation of bacteria to phage attack is usually caused by mutations on genes encoding phage receptors, which include lipopolysaccharides (LPS), outer membrane proteins, capsules, flagella, and pili, among others. However, some studies suggest that when phage resistance is caused by mutations in phage receptors, this might result in fitness costs to the resistance bacterium, which will ultimately become less virulent. Moreover, it has been shown that the evolution of bacterial resistance to phage attack changes the efflux pump mechanism, causing increased sensitivity to drugs from several antibiotic classes. Therefore, it is conceivable to think that phage therapy that uses phages that exert selection for multidrug-resistant bacteria to become antibiotic-sensitive could potentially reduce the incidence of antibiotic-resistant infections. Besides the prevention of phage adsorption by loss or modification of bacterial receptors, phage insensitivity can be caused by: prevention of phage DNA entry by superinfection exclusion systems; or degradation of phage DNA by restriction-modification systems or by CRISPR-Cas systems; and use of abortive infection systems that block phage replication, transcription, or translation, usually in conjunction with suicide of the host cell. Altogether, these mechanisms promote a quick adaptation of bacteria to phage attack and therefore, the emergence of phage-resistant mutants is frequent and unavoidable. It is still unclear whether the wide use of phages would cause resistance similar to what has been observed for antibiotics. In theory, this is not very likely to occur, since phages are very specific, and therefore, their selective pressure would affect a very narrow group of bacteria. However, we should also consider the fact that many phage resistance systems are mounted on mobile genetic elements, including prophages and plasmids, and thus may spread quite rapidly even without direct selection. Nevertheless, in contrast to antibiotics, phage preparations for therapeutic applications are expected to be developed in a personalized way because of the high specificity of phages. In addition, strategies have been proposed to counter the problem of phage resistance. One of the strategies is the use of phage cocktails with complementary host ranges (different host ranges, which, when combined, result in an overall broader host range) and targeting different bacterial receptors. Another strategy is the combination of phages with other antimicrobials such as antibiotics, disinfectants, or enzymes that could enhance their antibacterial activity. The genetic manipulation of phage genomes can also be a strategy to circumvent phage resistance. Safety aspects Bacteriophages are bacterial viruses, evolved to infect bacterial cells. To do that, phages must use characteristic structures at cell surfaces (receptors), and to propagate they need appropriate molecular tools inside the cells. Bacteria are prokaryotes, and their cells differ substantially from eukaryotes, including humans or animals. For this reason, phages meet the major safety requirement: they do not infect treated individuals. Even engineered phages and induced artificial internalization of phages into mammalian cells do not result in phage propagation. Natural transcytosis of unmodified phages, that is, uptake and internal transport to the other side of a cell, which was observed in human epithelial cells, did not result in phage propagation or cell damage. Recently, however, it was reported that filamentous temperate phages of P. aeruginosa can be endocytosed into human and murine leukocytes, resulting in transcription of the phage DNA. In turn, the product RNA triggers maladaptive innate viral pattern-recognition responses and thus inhibits the immune clearance of the bacteria. Whether this also applies to dsDNA phages like Caudovirales has not yet been established; this is an important question to be addressed as it may affect the overall safety of phage therapy. Due to many experimental treatments in human patients conducted in past decades, and to already existing RCTs (see section: Clinical experience and randomized controlled trials), phage safety can be assessed directly. The first safety trial in healthy human volunteers for a phage was conducted by Bruttin and Brüssow in 2005. They investigated the oral administration of Escherichia coli phage T4 and found no adverse effects of the treatment. Historical record shows that phages are safe, with mild side effects, if any. Still, administering bacteriophages can induce an immune response. Macrophages, key cells of the innate immune system, play a central role in mediating this response. The most frequent (though still rare) adverse reactions to phage preparations found in patients were symptoms from the digestive tract, local reactions at the site of administration of a phage preparation, superinfections, and a rise in body temperature. These reactions might have occurred because either toxins were released from bacteria destroyed by the phages—such toxin release from bacteria can also happen with antibiotic use—or due to leftover bacterial fragments or residual components from the bacterial growth medium ("food for bacteria") present in the phage treatment when unpurified preparations were used. When bacteriophages are introduced into the body, they may be recognized as foreign entities by macrophages through pattern recognition receptors (PRRs) such as Toll-like receptors (TLRs). The binding of bacteriophages to these receptors triggers macrophage activation, leading to phagocytosis (macrophages engulf and digest the bacteriophages) and cytokine production: activated macrophages produce pro-inflammatory cytokines. These cytokines can modulate the immune response but generally do not result in significant fever when phages are used appropriately. The route by which bacteriophages enter the body can affect the degree of immune activation. Applying bacteriophages directly to the mucosa targets the site of infection with minimal systemic exposure, leading to a localized immune response. Injecting bacteriophages into muscle tissue introduces them to a larger number of macrophages in the muscle and regional lymph nodes. In intravenous injection, direct introduction into the bloodstream exposes bacteriophages to macrophages throughout the body, including those in the spleen and liver. However, significant elevations in body temperature are uncommon and typically only observed in cases of rapid phage administration or high doses. Anticipating immune responses allows healthcare professionals to monitor patients appropriately and make treatment adjustments if necessary. Macrophages are integral to the body's immune response to bacteriophage therapy, mediating any potential immune reactions. Intravenous administration of bacteriophages is conducted under strict medical supervision, by specialists in infectious diseases within a hospital setting, due to potential adverse reactions. Adverse reactions to intravenous bacteriophage therapy may include hypotension, i.e., a drop in blood pressure, leading to loss of consciousness. A sudden drop (chills) and rise (fever) in body temperature, known as the Jarisch–Herxheimer reaction, can occur due to the rapid lysis of bacteria and release of endotoxins. Rapid bacterial lysis releases endotoxins (e.g., lipopolysaccharides from gram-negative bacteria) that trigger systemic inflammatory responses, including "cytokine storms". Continuous monitoring of heart rate, blood pressure, and temperature to detect early signs of adverse reactions is done after the intravenous phage administration. Successful treatment of life-threatening infections with intravenous phage therapy has been documented. Patients have responded to therapy after one or several intravenous administrations, clearing infections that were unresponsive to conventional treatments: phages can disrupt biofilms, which are often resistant to antibiotics, enhancing infection clearance. Bacteriophages must be produced in bacteria that are lysed (i.e., fragmented) during phage propagation. As such, phage lysates contain bacterial debris that may affect the human organism even when the phage itself is harmless. For these and other reasons, purification of bacteriophages is considered important, and phage preparations need to be assessed for their safety as a whole, particularly when phages are to be administered intravenously. This is consistent with general procedures for other drug candidates. In 2015, a group of phage therapy experts summarized the quality and safety requirements for sustainable phage therapy. Phage effects on the human microbiome also contribute to safety issues in phage therapy. Many phages, especially temperate ones, carry genes that can affect the pathogenicity of the host. Even lambda, a temperate phage of the E. coli K-12 laboratory strain, carries two genes that provide potential virulence benefits to the lysogenic host, one that increases intestinal adherence and the other that confers resistance to complement killing in the blood. For this reason, temperate phages are generally to be avoided as candidates for phage therapy, although in some cases, the lack of lytic phage candidates and emergency conditions may make such considerations moot. Another potential problem is generalized transduction, a term for the ability of some phages to transfer bacterial DNA from one host to another. This occurs because the systems for packaging of the phage DNA into capsids can mistakenly package host DNA instead. Indeed, with some well-characterized phages, up to 5% of the virus particles contain only bacterial DNA. Thus in a typical lysate, the entire genome of the propagating host is present in more than a million copies in every milliliter. For these reasons, it is imperative that any phage to be considered for therapeutic usage should be subjected to thorough genomic analysis and tested for the capacity for generalized transduction. As antibacterials, phages may also affect the composition of microbiomes, by infecting and killing phage-sensitive strains of bacteria. However, a major advantage of bacteriophages over antibiotics is the high specificity of bacteriophages. This specificity limits antibacterial activity to a sub-species level; typically, a phage kills only selected bacterial strains. For this reason, phages are much less likely (than antibiotics) to disturb the composition of a natural microbiome or to induce dysbiosis. This was demonstrated in experimental studies where microbiome composition was assessed by next-generation sequencing that revealed no important changes correlated with phage treatment in human treatments. Much of the difficulty in obtaining regulatory approval is proving to be the risks of using a self-replicating entity that has the capability to evolve. As with antibiotic therapy and other methods of countering bacterial infections, endotoxins are released by the bacteria as they are destroyed within the patient (Jarisch–Herxheimer reaction). This can cause symptoms of fever; in extreme cases, toxic shock (a problem also seen with antibiotics) is possible. Janakiraman Ramachandran argues that this complication can be avoided in those types of infection where this reaction is likely to occur by using genetically engineered bacteriophages that have had their gene responsible for producing endolysin removed. Without this gene, the host bacterium still dies but remains intact, because the lysis is disabled. On the other hand, this modification stops the exponential growth of phages, so one administered phage means at most one dead bacterial cell. Eventually, these dead cells are consumed by the normal house-cleaning duties of the phagocytes, which utilize enzymes to break down the whole bacterium and its contents into harmless proteins, polysaccharides, and lipids. Temperate (or lysogenic) bacteriophages are not generally used therapeutically, since this group can act as a way for bacteria to exchange DNA. This can help spread antibiotic resistance or even, theoretically, make the bacteria pathogenic, such as in cases of cholera. Carl Merril has claimed that harmless strains of corynebacterium may have been converted into C. diphtheriae that "probably killed a third of all Europeans who came to North America in the seventeenth century". Fortunately, many phages seem to be lytic only with negligible probability of becoming lysogenic. Regulation and legislation Approval of phage therapy for use in humans has not been given in Western countries, with a few exceptions. In the United States, Washington and Oregon law allows naturopathic physicians to use any therapy that is legal anywhere in the world on an experimental basis, and in Texas, phages are considered natural substances and can be used in addition to (but not as a replacement for) traditional therapy (they have been used routinely in a wound care clinic in Lubbock since 2006). In 2013, "the 20th biennial Evergreen International Phage Meeting ... conference drew 170 participants from 35 countries, including leaders of companies and institutes involved with human phage therapies from France, Australia, Georgia, Poland, and the United States." In France, phage therapy disappeared officially with the withdrawal of the Vidal dictionary (France's official drug directory), in 1978. The last phage preparation, produced by l'Institut du Bactériophage, was an ointment against skin infections. Phage therapy research ceased at about the same time across the country, with the closure of the bacteriophage department at the Pasteur Institute. Some hospital physicians continued to offer phage therapy until the 1990s, when production died out. On their rediscovery, at the end of the 1990s, phage preparations were classified as medicines, i.e., "medicinal products" in the EU or "drugs" in the US. However, the pharmaceutical legislation that had been implemented since their disappearance from Western medicine was mainly designed to cater for industrially-made pharmaceuticals, devoid of any customization and intended for large-scale distribution, and it was not deemed necessary to provide phage-specific requirements or concessions. Today's phage therapy products need to comply with the entire battery of medicinal product licensing requirements: manufacturing according to GMP, preclinical studies, phase I, II, and III clinical trials, and marketing authorisation. Technically, industrially produced predefined phage preparations could make it through the conventional pharmaceutical processes, minding some adaptations. However, phage specificity and resistance issues are likely to cause these defined preparations to have a relatively short useful lifespan. The pharmaceutical industry is currently not considering phage therapy products. Yet, a handful of small and medium-sized enterprises have shown interest, with the help of risk capital and/or public funding. Currently, no defined therapeutic phage product has made it to the EU or US markets. According to Jean-Paul Pirnay, therapeutic phages should be prepared individually and kept in large phage banks, ready to be used, upon testing for effectiveness against the patient's bacterial pathogen(s). Intermediary or combined (industrially made as well as precision phage preparations) approaches could be appropriate. However, it turns out to be difficult to reconcile classical phage therapy concepts, which are based on the timely adaptation of phage preparations, with current Western pharmaceutical R&D and marketing models. Repeated calls for a specific regulatory framework have not been heeded by European policymakers. A phage therapy framework based on the Biological Master File concept has been proposed as a (European) solution to regulatory issues, but European regulations do not allow for an extension of this concept to biologically active substances such as phages. Meanwhile, representatives from the medical, academic, and regulatory communities have established some (temporary) national solutions. For instance, phage applications have been performed in Europe under the umbrella of Article 37 (Unproven Interventions in Clinical Practice) of the Helsinki Declaration. To enable the application of phage therapy after Poland had joined the EU in 2004, the Ludwik Hirszfeld Institute of Immunology and Experimental Therapy in Wrocław opened its own Phage Therapy Unit (PTU). Phage therapy performed at the PTU is considered an "experimental treatment", covered by the adapted Act of 5 December 1996 on the Medical Profession (Polish Law Gazette, 2011, No. 277 item 1634) and Article 37 of the Helsinki Declaration. Similarly, in the last few years, a number of phage therapy interventions have been performed in the US under the FDA's emergency Investigational New Drug (eIND) protocol. Some patients have been treated with phages under the umbrella of "compassionate use", which is a treatment option that allows a physician to use a not-yet-authorized medicine in desperate cases. Under strict conditions, medicines under development can be made available for use in patients for whom no satisfactory authorized therapies are available and who cannot participate in clinical trials. In principle, this approach can only be applied to products for which earlier study results have demonstrated efficacy and safety, but have not yet been approved. Much like Article 37 of the Helsinki Declaration, the compassionate use treatment option can only be applied when the phages are expected to help in life-threatening or chronic and/or seriously debilitating diseases that are not treatable with formally approved products. In France, ANSM, the French medicine agency, has organized a specific committee—Comité Scientifique Spécialisé Temporaire (CSST)—for phage therapy, which consists of experts in various fields. Their task is to evaluate and guide each phage therapy request that ends up at the ANSM. Phage therapy requests are discussed together with the treating physicians and consensus advice is sent to the ANSM], which then decides whether or not to grant permission. Between 2006 and 2018, fifteen patients were treated in France (eleven recovered) using this pathway. In Belgium, in 2016 and in response to a number of parliamentary questions, Maggie De Block, the Minister of Social Affairs and Health, acknowledged that it is indeed not evident to treat phages as industrially made drugs, and therefore she proposed to investigate if the magistral preparation pathway could offer a solution. Magistral preparations (compounding pharmacies in the US) are not subjected to certain constraints such as GMP compliance and marketing authorization. As the "magistral preparation framework" was created to allow for adapted patient treatments and/or to use medicines for which there is no commercial interest, it seemed a suitable framework for precision phage therapy concepts. Magistral preparations are medicines prepared in a pharmacy in accordance with a medical prescription for an individual patient. They are made by a pharmacist (or under his/her supervision) from their constituent ingredients, according to the technical and scientific standards of pharmaceutical technology. Phage active pharmaceutical ingredients to be included in magistral preparations must meet the requirements of a monograph, which describes their production and quality control testing. They must be accompanied by a certificate of analysis, issued by a "Belgian Approved Laboratory", which has been granted an accreditation to perform batch-release testing of medicinal products. Since 2019, phages have been delivered in the form of magistral preparations to nominal patients in Belgium. The first phage therapy case in China can be traced back to 1958, at Shanghai Jiao Tong University School of Medicine. However, many regulations were not yet established back then, and phage therapy soon lost people's interest due to the prevalence of antibiotics, which eventually led to the antimicrobial resistance crisis. This prompted researchers in China as well as the Chinese government to pay attention to phage therapy again, and following the first investigator-initiated trial (IIT) by the Shanghai Institute of Phage in 2019, phage therapy rapidly flourished. Currently, commercial phage therapy applications must go through either one of two pathways. The first is for fixed-ingredient phage products. The second pathway is for personalized phage products, which need to go through IITs. This way, the products are considered restrictive medical technologies. Application in other species Animals Phage therapy has been a relevant mode of treatment in animals for decades. It has been proposed as a method of treating bacterial infections in the veterinary medical field in response to the rampant use of antibiotics. Studies have investigated the application of phage therapy in livestock species as well as companion animals. Brigham Young University has been researching the use of phage therapy to treat American foulbrood in honeybees. Phage therapy is also being investigated for potential applications in aquaculture. Plants Phage therapy has been studied for bacterial spot of stonefruit, caused by Xanthomonas pruni (syn. X. campestris pv. pruni, syn. X. arboricola pv. pruni) in prunus species. Some treatments have been very successful. Cultural impact The 1925 novel and 1926 Pulitzer Prize winner Arrowsmith by Sinclair Lewis used phage therapy as a plot point. Greg Bear's 2002 novel Vitals features phage therapy, based on Soviet research, used to transfer genetic material. The 2012 collection of military history essays about the changing role of women in warfare, Women in War – From Home Front to Front Line includes a chapter featuring phage therapy: "Chapter 17: Women who thawed the Cold War". Steffanie A. Strathdee's book The Perfect Predator: An Epidemiologist's Journey to Save Her Husband from a Deadly Superbug, co-written with her husband, Thomas Patterson, was published by Hachette Book Group in 2019. It describes Strathdee's ultimately successful attempt to introduce phage therapy as a life-saving treatment for her husband, critically ill with a completely antibiotic-resistant Acinetobacter baumannii infection following severe pancreatitis.	Biology and health sciences	Treatments	Health
682478	https://en.wikipedia.org/wiki/Chinese%20cabbage	Chinese cabbage	Chinese cabbage (Brassica rapa, subspecies pekinensis and chinensis) is either of two cultivar groups of leaf vegetables often used in Chinese cuisine: the Pekinensis Group (napa cabbage) and the Chinensis Group (bok choy). These vegetables are both variant cultivars or subspecies of B. rapa and belong to the same genus as such Western staples as cabbage, broccoli, and cauliflower. Both have many variations in name, spelling, and scientific classification, especially bok choy cultivars. History The Chinese cabbage was principally grown in the Yangtze River Delta region, but the Ming dynasty naturalist Li Shizhen popularized it by bringing attention to its medicinal qualities. The variant cultivated in Zhejiang around the 14th century was brought north, and the northern harvest of napa cabbage soon exceeded the southern one. These were then exported back south along the Grand Canal to Hangzhou and traded by sea as far south as Guangdong. Napa cabbage became a staple in Northeastern Chinese cuisine for making suan cai, Chinese sauerkraut. In Korea, napa cabbage became baek-kimchi, which developed into kimchi. Chinese cabbage is now commonly found in markets throughout the world, catering both to the Chinese diaspora and to northern markets that appreciate its resistance to cold. In 2017, aboard the International Space Station, a crop of Chinese cabbage from a plant growth device included an allotment for crew consumption, while the rest was saved for scientific study. Cultivar groups There are two distinct groups of Brassica rapa used as leaf vegetables in China, and a wide range of cultivars within these two groups. The binomial name B. campestris is also used. Pekinensis Group This group is the more common of the two, especially outside Asia; names such as napa cabbage, dà báicài (, "large white vegetable"); Baguio petsay or petsay wombok (Tagalog); Chinese white cabbage; "wong a pak" (Hokkien, Fujianese); baechu (), wongbok; hakusai () and "suann-tang-pe̍h-á" (Taiwanese) usually refer to members of this group. Pekinensis Group cabbages have broad green leaves with white petioles, tightly wrapped in a cylindrical formation and usually forming a compact head. As the group name indicates, this is particularly popular in northern China around Beijing (Peking). Chinensis Group Chinensis Group cultivars do not form heads; instead, they have smooth, dark green leaf blades forming a cluster reminiscent of mustard or celery. These cultivars are popular in southern China and Southeast Asia. Being winter-hardy, they are increasingly grown in Northern Europe. This group was originally classified as its own species under the name B. chinensis by Linnaeus.	Biology and health sciences	Brassicales	null
682482	https://en.wikipedia.org/wiki/Human	Human	Humans (Homo sapiens) or modern humans are the most common and widespread species of primate, and the last surviving species of the genus Homo. They are great apes characterized by their hairlessness, bipedalism, and high intelligence. Humans have large brains, enabling more advanced cognitive skills that enable them to thrive and adapt in varied environments, develop highly complex tools, and form complex social structures and civilizations. Humans are highly social, with individual humans tending to belong to a multi-layered network of distinct social groupsfrom families and peer groups to corporations and political states. As such, social interactions between humans have established a wide variety of values, social norms, languages, and traditions (collectively termed institutions), each of which bolsters human society. Humans are also highly curious: the desire to understand and influence phenomena has motivated humanity's development of science, technology, philosophy, mythology, religion, and other frameworks of knowledge; humans also study themselves through such domains as anthropology, social science, history, psychology, and medicine. As of January 2025, there are estimated to be more than 8 billion humans alive. For most of their history, humans were nomadic hunter-gatherers. Humans began exhibiting behavioral modernity about 160,000–60,000 years ago. The Neolithic Revolution, which began in Southwest Asia around 13,000 years ago (and separately in a few other places), saw the emergence of agriculture and permanent human settlement; in turn, this led to the development of civilization and kickstarted a period of continuous (and ongoing) population growth and rapid technological change. Since then, a number of civilizations have risen and fallen, while a number of sociocultural and technological developments have resulted in significant changes to the human lifestyle. Although some scientists equate the term "humans" with all members of the genus Homo, in common usage it generally refers to Homo sapiens, the only extant member. All other members of the genus Homo, which are now extinct, are known as archaic humans, and the term "modern human" is used to distinguish Homo sapiens from archaic humans. Anatomically modern humans emerged around 300,000 years ago in Africa, evolving from Homo heidelbergensis or a similar species. Migrating out of Africa, they gradually replaced and interbred with local populations of archaic humans. Multiple hypotheses for the extinction of archaic human species such as Neanderthals include competition, violence, interbreeding with Homo sapiens, or inability to adapt to climate change. Genes and the environment influence human biological variation in visible characteristics, physiology, disease susceptibility, mental abilities, body size, and life span. Though humans vary in many traits (such as genetic predispositions and physical features), humans are among the least genetically diverse primates. Any two humans are at least 99% genetically similar. Humans are sexually dimorphic: generally, males have greater body strength and females have a higher body fat percentage. At puberty, humans develop secondary sex characteristics. Females are capable of pregnancy, usually between puberty, at around 12 years old, and menopause, around the age of 50. Humans are omnivorous, capable of consuming a wide variety of plant and animal material, and have used fire and other forms of heat to prepare and cook food since the time of Homo erectus. Humans have had a dramatic effect on the environment. They are apex predators, being rarely preyed upon by other species. Human population growth, industrialization, land development, overconsumption and combustion of fossil fuels have led to environmental destruction and pollution that significantly contributes to the ongoing mass extinction of other forms of life. Within the last century, humans have explored challenging environments such as Antarctica, the deep sea, and outer space. Human habitation within these hostile environments is restrictive and expensive, typically limited in duration, and restricted to scientific, military, or industrial expeditions. Humans have visited the Moon and made their presence known on other celestial bodies through human-made robotic spacecraft. Since the early 20th century, there has been continuous human presence in Antarctica through research stations and, since 2000, in space through habitation on the International Space Station. Humans can survive for up to eight weeks without food and several days without water. Humans are generally diurnal, sleeping on average seven to nine hours per day. Childbirth is dangerous, with a high risk of complications and death. Often, both the mother and the father provide care for their children, who are helpless at birth. Etymology and definition All modern humans are classified into the species Homo sapiens, coined by Carl Linnaeus in his 1735 work Systema Naturae. The generic name Homo is a learned 18th-century derivation from Latin , which refers to humans of either sex. The word human can refer to all members of the Homo genus. The name Homo sapiens means 'wise man' or 'knowledgeable man'. There is disagreement if certain extinct members of the genus, namely Neanderthals, should be included as a separate species of humans or as a subspecies of H. sapiens. Human is a loanword of Middle English from Old French , ultimately from Latin , the adjectival form of ('man'in the sense of humanity). The native English term man can refer to the species generally (a synonym for humanity) as well as to human males. It may also refer to individuals of either sex. Despite the fact that the word animal is colloquially used as an antonym for human, and contrary to a common biological misconception, humans are animals. The word person is often used interchangeably with human, but philosophical debate exists as to whether personhood applies to all humans or all sentient beings, and further if a human can lose personhood (such as by going into a persistent vegetative state). Evolution Humans are apes (superfamily Hominoidea). The lineage of apes that eventually gave rise to humans first split from gibbons (family Hylobatidae) and orangutans (genus Pongo), then gorillas (genus Gorilla), and finally, chimpanzees and bonobos (genus Pan). The last split, between the human and chimpanzee–bonobo lineages, took place around 8–4 million years ago, in the late Miocene epoch. During this split, chromosome 2 was formed from the joining of two other chromosomes, leaving humans with only 23 pairs of chromosomes, compared to 24 for the other apes. Following their split with chimpanzees and bonobos, the hominins diversified into many species and at least two distinct genera. All but one of these lineagesrepresenting the genus Homo and its sole extant species Homo sapiensare now extinct. The genus Homo evolved from Australopithecus. Though fossils from the transition are scarce, the earliest members of Homo share several key traits with Australopithecus. The earliest record of Homo is the 2.8 million-year-old specimen LD 350-1 from Ethiopia, and the earliest named species are Homo habilis and Homo rudolfensis which evolved by 2.3 million years ago. H. erectus (the African variant is sometimes called H. ergaster) evolved 2 million years ago and was the first archaic human species to leave Africa and disperse across Eurasia. H. erectus also was the first to evolve a characteristically human body plan. Homo sapiens emerged in Africa around 300,000 years ago from a species commonly designated as either H. heidelbergensis or H. rhodesiensis, the descendants of H. erectus that remained in Africa. H. sapiens migrated out of the continent, gradually replacing or interbreeding with local populations of archaic humans. Humans began exhibiting behavioral modernity about 160,000–70,000 years ago, and possibly earlier. This development was likely selected amidst natural climate change in Middle to Late Pleistocene Africa. The "out of Africa" migration took place in at least two waves, the first around 130,000 to 100,000 years ago, the second (Southern Dispersal) around 70,000 to 50,000 years ago. H. sapiens proceeded to colonize all the continents and larger islands, arriving in Eurasia 125,000 years ago, Australia around 65,000 years ago, the Americas around 15,000 years ago, and remote islands such as Hawaii, Easter Island, Madagascar, and New Zealand in the years 300 to 1280 CE. Human evolution was not a simple linear or branched progression but involved interbreeding between related species. Genomic research has shown that hybridization between substantially diverged lineages was common in human evolution. DNA evidence suggests that several genes of Neanderthal origin are present among all non sub-Saharan-African populations, and Neanderthals and other hominins, such as Denisovans, may have contributed up to 6% of their genome to present-day non sub-Saharan-African humans. Human evolution is characterized by a number of morphological, developmental, physiological, and behavioral changes that have taken place since the split between the last common ancestor of humans and chimpanzees. The most significant of these adaptations are hairlessness, obligate bipedalism, increased brain size and decreased sexual dimorphism (neoteny). The relationship between all these changes is the subject of ongoing debate. History Prehistory Until about 12,000 years ago, all humans lived as hunter-gatherers. The Neolithic Revolution (the invention of agriculture) first took place in Southwest Asia and spread through large parts of the Old World over the following millennia. It also occurred independently in Mesoamerica (about 6,000 years ago), China, Papua New Guinea, and the Sahel and West Savanna regions of Africa. Access to food surplus led to the formation of permanent human settlements, the domestication of animals and the use of metal tools for the first time in history. Agriculture and sedentary lifestyle led to the emergence of early civilizations. Ancient An urban revolution took place in the 4th millennium BCE with the development of city-states, particularly Sumerian cities located in Mesopotamia. It was in these cities that the earliest known form of writing, cuneiform script, appeared around 3000 BCE. Other major civilizations to develop around this time were Ancient Egypt and the Indus Valley Civilisation. They eventually traded with each other and invented technology such as wheels, plows and sails. Emerging by 3000 BCE, the Caral–Supe civilization is the oldest complex civilization in the Americas. Astronomy and mathematics were also developed and the Great Pyramid of Giza was built. There is evidence of a severe drought lasting about a hundred years that may have caused the decline of these civilizations, with new ones appearing in the aftermath. Babylonians came to dominate Mesopotamia while others, such as the Poverty Point culture, Minoans and the Shang dynasty, rose to prominence in new areas. The Late Bronze Age collapse around 1200 BCE resulted in the disappearance of a number of civilizations and the beginning of the Greek Dark Ages. During this period iron started replacing bronze, leading to the Iron Age. In the 5th century BCE, history started being recorded as a discipline, which provided a much clearer picture of life at the time. Between the 8th and 6th century BCE, Europe entered the classical antiquity age, a period when ancient Greece and ancient Rome flourished. Around this time other civilizations also came to prominence. The Maya civilization started to build cities and create complex calendars. In Africa, the Kingdom of Aksum overtook the declining Kingdom of Kush and facilitated trade between India and the Mediterranean. In West Asia, the Achaemenid Empire's system of centralized governance became the precursor to many later empires, while the Gupta Empire in India and the Han dynasty in China have been described as golden ages in their respective regions. Medieval Following the fall of the Western Roman Empire in 476, Europe entered the Middle Ages. During this period, Christianity and the Church would provide centralized authority and education. In the Middle East, Islam became the prominent religion and expanded into North Africa. It led to an Islamic Golden Age, inspiring achievements in architecture, the revival of old advances in science and technology, and the formation of a distinct way of life. The Christian and Islamic worlds would eventually clash, with the Kingdom of England, the Kingdom of France and the Holy Roman Empire declaring a series of holy wars to regain control of the Holy Land from Muslims. In the Americas, between 200 and 900 CE Mesoamerica was in its Classic Period, while further north, complex Mississippian societies would arise starting around 800 CE. The Mongol Empire would conquer much of Eurasia in the 13th and 14th centuries. Over this same time period, the Mali Empire in Africa grew to be the largest empire on the continent, stretching from Senegambia to Ivory Coast. Oceania would see the rise of the Tuʻi Tonga Empire which expanded across many islands in the South Pacific. By the late 15th century, the Aztecs and Inca had become the dominant power in Mesoamerica and the Andes, respectively. Modern The early modern period in Europe and the Near East (–1800) began with the final defeat of the Byzantine Empire, and the rise of the Ottoman Empire. Meanwhile, Japan entered the Edo period, the Qing dynasty rose in China and the Mughal Empire ruled much of India. Europe underwent the Renaissance, starting in the 15th century, and the Age of Discovery began with the exploring and colonizing of new regions. This included the colonization of the Americas and the Columbian Exchange. This expansion led to the Atlantic slave trade and the genocide of Native American peoples. This period also marked the Scientific Revolution, with great advances in mathematics, mechanics, astronomy and physiology. The late modern period (1800–present) saw the Technological and Industrial Revolution bring such discoveries as imaging technology, major innovations in transport and energy development. Influenced by Enlightenment ideals, the Americas and Europe experienced a period of political revolutions known as the Age of Revolution. The Napoleonic Wars raged through Europe in the early 1800s, Spain lost most of its colonies in the New World, while Europeans continued expansion into Africawhere European control went from 10% to almost 90% in less than 50 yearsand Oceania. In the 19th century, the British Empire expanded to become the world's largest empire.A tenuous balance of power among European nations collapsed in 1914 with the outbreak of the First World War, one of the deadliest conflicts in history. In the 1930s, a worldwide economic crisis led to the rise of authoritarian regimes and a Second World War, involving almost all of the world's countries. The war's destruction led to the collapse of most global empires, leading to widespread decolonization. Following the conclusion of the Second World War in 1945, the United States and the USSR emerged as the remaining global superpowers. This led to a Cold War that saw a struggle for global influence, including a nuclear arms race and a space race, ending in the collapse of the Soviet Union. The current Information Age, spurred by the development of the Internet and artificial intelligence systems, sees the world becoming increasingly globalized and interconnected. Habitat and population Early human settlements were dependent on proximity to water anddepending on the lifestyleother natural resources used for subsistence, such as populations of animal prey for hunting and arable land for growing crops and grazing livestock. Modern humans, however, have a great capacity for altering their habitats by means of technology, irrigation, urban planning, construction, deforestation and desertification. Human settlements continue to be vulnerable to natural disasters, especially those placed in hazardous locations and with low quality of construction. Grouping and deliberate habitat alteration is often done with the goals of providing protection, accumulating comforts or material wealth, expanding the available food, improving aesthetics, increasing knowledge or enhancing the exchange of resources. Humans are one of the most adaptable species, despite having a low or narrow tolerance for many of the earth's extreme environments. Currently the species is present in all eight biogeographical realms, although their presence in the Antarctic realm is very limited to research stations and annually there is a population decline in the winter months of this realm. Humans established nation-states in the other seven realms, such as South Africa, India, Russia, Australia, Fiji, United States and Brazil (each located in a different biogeographical realm). By using advanced tools and clothing, humans have been able to extend their tolerance to a wide variety of temperatures, humidities, and altitudes. As a result, humans are a cosmopolitan species found in almost all regions of the world, including tropical rainforest, arid desert, extremely cold arctic regions, and heavily polluted cities; in comparison, most other species are confined to a few geographical areas by their limited adaptability. The human population is not, however, uniformly distributed on the Earth's surface, because the population density varies from one region to another, and large stretches of surface are almost completely uninhabited, like Antarctica and vast swathes of the ocean. Most humans (61%) live in Asia; the remainder live in the Americas (14%), Africa (14%), Europe (11%), and Oceania (0.5%). Estimates of the population at the time agriculture emerged in around 10,000 BC have ranged between 1 million and 15 million. Around 50–60 million people lived in the combined eastern and western Roman Empire in the 4th century AD. Bubonic plagues, first recorded in the 6th century AD, reduced the population by 50%, with the Black Death killing 75–200 million people in Eurasia and North Africa alone. Human population is believed to have reached one billion in 1800. It has since then increased exponentially, reaching two billion in 1930 and three billion in 1960, four in 1975, five in 1987 and six billion in 1999. It passed seven billion in 2011 and passed eight billion in November 2022. It took over two million years of human prehistory and history for the human population to reach one billion and only 207 years more to grow to 7 billion. The combined biomass of the carbon of all the humans on Earth in 2018 was estimated at 60 million tons, about 10 times larger than that of all non-domesticated mammals. In 2018, 4.2 billion humans (55%) lived in urban areas, up from 751 million in 1950. The most urbanized regions are Northern America (82%), Latin America (81%), Europe (74%) and Oceania (68%), with Africa and Asia having nearly 90% of the world's 3.4 billion rural population. Problems for humans living in cities include various forms of pollution and crime, especially in inner city and suburban slums. Biology Anatomy and physiology Most aspects of human physiology are closely homologous to corresponding aspects of animal physiology. The dental formula of humans is: . Humans have proportionately shorter palates and much smaller teeth than other primates. They are the only primates to have short, relatively flush canine teeth. Humans have characteristically crowded teeth, with gaps from lost teeth usually closing up quickly in young individuals. Humans are gradually losing their third molars, with some individuals having them congenitally absent. Humans share with chimpanzees a vestigial tail, appendix, flexible shoulder joints, grasping fingers and opposable thumbs. Humans also have a more barrel-shaped chests in contrast to the funnel shape of other apes, an adaptation for bipedal respiration. Apart from bipedalism and brain size, humans differ from chimpanzees mostly in smelling, hearing and digesting proteins. While humans have a density of hair follicles comparable to other apes, it is predominantly vellus hair, most of which is so short and wispy as to be practically invisible. Humans have about 2 million sweat glands spread over their entire bodies, many more than chimpanzees, whose sweat glands are scarce and are mainly located on the palm of the hand and on the soles of the feet. It is estimated that the worldwide average height for an adult human male is about , while the worldwide average height for adult human females is about . Shrinkage of stature may begin in middle age in some individuals but tends to be typical in the extremely aged. Throughout history, human populations have universally become taller, probably as a consequence of better nutrition, healthcare, and living conditions. The average mass of an adult human is for females and for males. Like many other conditions, body weight and body type are influenced by both genetic susceptibility and environment and varies greatly among individuals. Humans have a far faster and more accurate throw than other animals. Humans are also among the best long-distance runners in the animal kingdom, but slower over short distances. Humans' thinner body hair and more productive sweat glands help avoid heat exhaustion while running for long distances. Compared to other apes, the human heart produces greater stroke volume and cardiac output and the aorta is proportionately larger. Genetics Like most animals, humans are a diploid and eukaryotic species. Each somatic cell has two sets of 23 chromosomes, each set received from one parent; gametes have only one set of chromosomes, which is a mixture of the two parental sets. Among the 23 pairs of chromosomes, there are 22 pairs of autosomes and one pair of sex chromosomes. Like other mammals, humans have an XY sex-determination system, so that females have the sex chromosomes XX and males have XY. Genes and environment influence human biological variation in visible characteristics, physiology, disease susceptibility and mental abilities. The exact influence of genes and environment on certain traits is not well understood. While no humansnot even monozygotic twinsare genetically identical, two humans on average will have a genetic similarity of 99.5%-99.9%. This makes them more homogeneous than other great apes, including chimpanzees. This small variation in human DNA compared to many other species suggests a population bottleneck during the Late Pleistocene (around 100,000 years ago), in which the human population was reduced to a small number of breeding pairs. The forces of natural selection have continued to operate on human populations, with evidence that certain regions of the genome display directional selection in the past 15,000 years. The human genome was first sequenced in 2001 and by 2020 hundreds of thousands of genomes had been sequenced. In 2012 the International HapMap Project had compared the genomes of 1,184 individuals from 11 populations and identified 1.6 million single nucleotide polymorphisms. African populations harbor the highest number of private genetic variants. While many of the common variants found in populations outside of Africa are also found on the African continent, there are still large numbers that are private to these regions, especially Oceania and the Americas. By 2010 estimates, humans have approximately 22,000 genes. By comparing mitochondrial DNA, which is inherited only from the mother, geneticists have concluded that the last female common ancestor whose genetic marker is found in all modern humans, the so-called mitochondrial Eve, must have lived around 90,000 to 200,000 years ago. Life cycle Most human reproduction takes place by internal fertilization via sexual intercourse, but can also occur through assisted reproductive technology procedures. The average gestation period is 38 weeks, but a normal pregnancy can vary by up to 37 days. Embryonic development in the human covers the first eight weeks of development; at the beginning of the ninth week the embryo is termed a fetus. Humans are able to induce early labor or perform a caesarean section if the child needs to be born earlier for medical reasons. In developed countries, infants are typically in weight and in height at birth. However, low birth weight is common in developing countries, and contributes to the high levels of infant mortality in these regions. Compared with other species, human childbirth is dangerous, with a much higher risk of complications and death. The size of the fetus's head is more closely matched to the pelvis than in other primates. The reason for this is not completely understood, but it contributes to a painful labor that can last 24 hours or more. The chances of a successful labor increased significantly during the 20th century in wealthier countries with the advent of new medical technologies. In contrast, pregnancy and natural childbirth remain hazardous ordeals in developing regions of the world, with maternal death rates approximately 100 times greater than in developed countries. Both the mother and the father provide care for human offspring, in contrast to other primates, where parental care is mostly done by the mother. Helpless at birth, humans continue to grow for some years, typically reaching sexual maturity at 15 to 17 years of age. The human life span has been split into various stages ranging from three to twelve. Common stages include infancy, childhood, adolescence, adulthood and old age. The lengths of these stages have varied across cultures and time periods but is typified by an unusually rapid growth spurt during adolescence. Human females undergo menopause and become infertile at around the age of 50. It has been proposed that menopause increases a woman's overall reproductive success by allowing her to invest more time and resources in her existing offspring, and in turn their children (the grandmother hypothesis), rather than by continuing to bear children into old age. The life span of an individual depends on two major factors, genetics and lifestyle choices. For various reasons, including biological/genetic causes, women live on average about four years longer than men. , the global average life expectancy at birth of a girl is estimated to be 74.9 years compared to 70.4 for a boy. There are significant geographical variations in human life expectancy, mostly correlated with economic developmentfor example, life expectancy at birth in Hong Kong is 87.6 years for girls and 81.8 for boys, while in the Central African Republic, it is 55.0 years for girls and 50.6 for boys. The developed world is generally aging, with the median age around 40 years. In the developing world, the median age is between 15 and 20 years. While one in five Europeans is 60 years of age or older, only one in twenty Africans is 60 years of age or older. In 2012, the United Nations estimated that there were 316,600 living centenarians (humans of age 100 or older) worldwide. Diet Humans are omnivorous, capable of consuming a wide variety of plant and animal material. Human groups have adopted a range of diets from purely vegan to primarily carnivorous. In some cases, dietary restrictions in humans can lead to deficiency diseases; however, stable human groups have adapted to many dietary patterns through both genetic specialization and cultural conventions to use nutritionally balanced food sources. The human diet is prominently reflected in human culture and has led to the development of food science. Until the development of agriculture, Homo sapiens employed a hunter-gatherer method as their sole means of food collection. This involved combining stationary food sources (such as fruits, grains, tubers, and mushrooms, insect larvae and aquatic mollusks) with wild game, which must be hunted and captured in order to be consumed. It has been proposed that humans have used fire to prepare and cook food since the time of Homo erectus. Human domestication of wild plants began about 11,700 years ago, leading to the development of agriculture, a gradual process called the Neolithic Revolution. These dietary changes may also have altered human biology; the spread of dairy farming provided a new and rich source of food, leading to the evolution of the ability to digest lactose in some adults. The types of food consumed, and how they are prepared, have varied widely by time, location, and culture. In general, humans can survive for up to eight weeks without food, depending on stored body fat. Survival without water is usually limited to three or four days, with a maximum of one week. In 2020 it is estimated 9 million humans die every year from causes directly or indirectly related to starvation. Childhood malnutrition is also common and contributes to the global burden of disease. However, global food distribution is not even, and obesity among some human populations has increased rapidly, leading to health complications and increased mortality in some developed and a few developing countries. Worldwide, over one billion people are obese, while in the United States 35% of people are obese, leading to this being described as an "obesity epidemic." Obesity is caused by consuming more calories than are expended, so excessive weight gain is usually caused by an energy-dense diet. Biological variation There is biological variation in the human specieswith traits such as blood type, genetic diseases, cranial features, facial features, organ systems, eye color, hair color and texture, height and build, and skin color varying across the globe. The typical height of an adult human is between , although this varies significantly depending on sex, ethnic origin, and family bloodlines. Body size is partly determined by genes and is also significantly influenced by environmental factors such as diet, exercise, and sleep patterns. There is evidence that populations have adapted genetically to various external factors. The genes that allow adult humans to digest lactose are present in high frequencies in populations that have long histories of cattle domestication and are more dependent on cow milk. Sickle cell anemia, which may provide increased resistance to malaria, is frequent in populations where malaria is endemic. Populations that have for a very long time inhabited specific climates tend to have developed specific phenotypes that are beneficial for those environmentsshort stature and stocky build in cold regions, tall and lanky in hot regions, and with high lung capacities or other adaptations at high altitudes. Some populations have evolved highly unique adaptations to very specific environmental conditions, such as those advantageous to ocean-dwelling lifestyles and freediving in the Bajau. Human hair ranges in color from red to blond to brown to black, which is the most frequent. Hair color depends on the amount of melanin, with concentrations fading with increased age, leading to grey or even white hair. Skin color can range from darkest brown to lightest peach, or even nearly white or colorless in cases of albinism. It tends to vary clinally and generally correlates with the level of ultraviolet radiation in a particular geographic area, with darker skin mostly around the equator. Skin darkening may have evolved as protection against ultraviolet solar radiation. Light skin pigmentation protects against depletion of vitamin D, which requires sunlight to make. Human skin also has a capacity to darken (tan) in response to exposure to ultraviolet radiation. There is relatively little variation between human geographical populations, and most of the variation that occurs is at the individual level. Much of human variation is continuous, often with no clear points of demarcation. Genetic data shows that no matter how population groups are defined, two people from the same population group are almost as different from each other as two people from any two different population groups. Dark-skinned populations that are found in Africa, Australia, and South Asia are not closely related to each other. Genetic research has demonstrated that human populations native to the African continent are the most genetically diverse and genetic diversity decreases with migratory distance from Africa, possibly the result of bottlenecks during human migration. These non-African populations acquired new genetic inputs from local admixture with archaic populations and have much greater variation from Neanderthals and Denisovans than is found in Africa, though Neanderthal admixture into African populations may be underestimated. Furthermore, recent studies have found that populations in sub-Saharan Africa, and particularly West Africa, have ancestral genetic variation which predates modern humans and has been lost in most non-African populations. Some of this ancestry is thought to originate from admixture with an unknown archaic hominin that diverged before the split of Neanderthals and modern humans. Humans are a gonochoric species, meaning they are divided into male and female sexes. The greatest degree of genetic variation exists between males and females. While the nucleotide genetic variation of individuals of the same sex across global populations is no greater than 0.1%–0.5%, the genetic difference between males and females is between 1% and 2%. Males on average are 15% heavier and taller than females. On average, men have about 40–50% more upper-body strength and 20–30% more lower-body strength than women at the same weight, due to higher amounts of muscle and larger muscle fibers. Women generally have a higher body fat percentage than men. Women have lighter skin than men of the same population; this has been explained by a higher need for vitamin D in females during pregnancy and lactation. As there are chromosomal differences between females and males, some X and Y chromosome-related conditions and disorders only affect either men or women. After allowing for body weight and volume, the male voice is usually an octave deeper than the female voice. Women have a longer life span in almost every population around the world. There are intersex conditions in the human population, however these are rare. Psychology The human brain, the focal point of the central nervous system in humans, controls the peripheral nervous system. In addition to controlling "lower", involuntary, or primarily autonomic activities such as respiration and digestion, it is also the locus of "higher" order functioning such as thought, reasoning, and abstraction. These cognitive processes constitute the mind, and, along with their behavioral consequences, are studied in the field of psychology. Humans have a larger and more developed prefrontal cortex than other primates, the region of the brain associated with higher cognition. This has led humans to proclaim themselves to be more intelligent than any other known species. Objectively defining intelligence is difficult, with other animals adapting senses and excelling in areas that humans are unable to. There are some traits that, although not strictly unique, do set humans apart from other animals. Humans may be the only animals who have episodic memory and who can engage in "mental time travel". Even compared with other social animals, humans have an unusually high degree of flexibility in their facial expressions. Humans are the only animals known to cry emotional tears. Humans are one of the few animals able to self-recognize in mirror tests and there is also debate over to what extent humans are the only animals with a theory of mind. Sleep and dreaming Humans are generally diurnal. The average sleep requirement is between seven and nine hours per day for an adult and nine to ten hours per day for a child; elderly people usually sleep for six to seven hours. Having less sleep than this is common among humans, even though sleep deprivation can have negative health effects. A sustained restriction of adult sleep to four hours per day has been shown to correlate with changes in physiology and mental state, including reduced memory, fatigue, aggression, and bodily discomfort. During sleep humans dream, where they experience sensory images and sounds. Dreaming is stimulated by the pons and mostly occurs during the REM phase of sleep. The length of a dream can vary, from a few seconds up to 30 minutes. Humans have three to five dreams per night, and some may have up to seven. Dreamers are more likely to remember the dream if awakened during the REM phase. The events in dreams are generally outside the control of the dreamer, with the exception of lucid dreaming, where the dreamer is self-aware. Dreams can at times make a creative thought occur or give a sense of inspiration. Consciousness and thought Human consciousness, at its simplest, is sentience or awareness of internal or external existence. Despite centuries of analyses, definitions, explanations and debates by philosophers and scientists, consciousness remains puzzling and controversial, being "at once the most familiar and most mysterious aspect of our lives". The only widely agreed notion about the topic is the intuition that it exists. Opinions differ about what exactly needs to be studied and explained as consciousness. Some philosophers divide consciousness into phenomenal consciousness, which is sensory experience itself, and access consciousness, which can be used for reasoning or directly controlling actions. It is sometimes synonymous with 'the mind', and at other times, an aspect of it. Historically it is associated with introspection, private thought, imagination and volition. It now often includes some kind of experience, cognition, feeling or perception. It may be 'awareness', or 'awareness of awareness', or self-awareness. There might be different levels or orders of consciousness, or different kinds of consciousness, or just one kind with different features. The process of acquiring knowledge and understanding through thought, experience, and the senses is known as cognition. The human brain perceives the external world through the senses, and each individual human is influenced greatly by his or her experiences, leading to subjective views of existence and the passage of time. The nature of thought is central to psychology and related fields. Cognitive psychology studies cognition, the mental processes underlying behavior. Largely focusing on the development of the human mind through the life span, developmental psychology seeks to understand how people come to perceive, understand, and act within the world and how these processes change as they age. This may focus on intellectual, cognitive, neural, social, or moral development. Psychologists have developed intelligence tests and the concept of intelligence quotient in order to assess the relative intelligence of human beings and study its distribution among population. Motivation and emotion Human motivation is not yet wholly understood. From a psychological perspective, Maslow's hierarchy of needs is a well-established theory that can be defined as the process of satisfying certain needs in ascending order of complexity. From a more general, philosophical perspective, human motivation can be defined as a commitment to, or withdrawal from, various goals requiring the application of human ability. Furthermore, incentive and preference are both factors, as are any perceived links between incentives and preferences. Volition may also be involved, in which case willpower is also a factor. Ideally, both motivation and volition ensure the selection, striving for, and realization of goals in an optimal manner, a function beginning in childhood and continuing throughout a lifetime in a process known as socialization. Emotions are biological states associated with the nervous system brought on by neurophysiological changes variously associated with thoughts, feelings, behavioral responses, and a degree of pleasure or displeasure. They are often intertwined with mood, temperament, personality, disposition, creativity, and motivation. Emotion has a significant influence on human behavior and their ability to learn. Acting on extreme or uncontrolled emotions can lead to social disorder and crime, with studies showing criminals may have a lower emotional intelligence than normal. Emotional experiences perceived as pleasant, such as joy, interest or contentment, contrast with those perceived as unpleasant, like anxiety, sadness, anger, and despair. Happiness, or the state of being happy, is a human emotional condition. The definition of happiness is a common philosophical topic. Some define it as experiencing the feeling of positive emotional affects, while avoiding the negative ones. Others see it as an appraisal of life satisfaction or quality of life. Recent research suggests that being happy might involve experiencing some negative emotions when humans feel they are warranted. Sexuality and love For humans, sexuality involves biological, erotic, physical, emotional, social, or spiritual feelings and behaviors. Because it is a broad term, which has varied with historical contexts over time, it lacks a precise definition. The biological and physical aspects of sexuality largely concern the human reproductive functions, including the human sexual response cycle. Sexuality also affects and is affected by cultural, political, legal, philosophical, moral, ethical, and religious aspects of life. Sexual desire, or libido, is a basic mental state present at the beginning of sexual behavior. Studies show that men desire sex more than women and masturbate more often. Humans can fall anywhere along a continuous scale of sexual orientation, although most humans are heterosexual. While homosexual behavior occurs in some other animals, only humans and domestic sheep have so far been found to exhibit exclusive preference for same-sex relationships. Most evidence supports nonsocial, biological causes of sexual orientation, as cultures that are very tolerant of homosexuality do not have significantly higher rates of it. Research in neuroscience and genetics suggests that other aspects of human sexuality are biologically influenced as well. Love most commonly refers to a feeling of strong attraction or emotional attachment. It can be impersonal (the love of an object, ideal, or strong political or spiritual connection) or interpersonal (love between humans). When in love dopamine, norepinephrine, serotonin and other chemicals stimulate the brain's pleasure center, leading to side effects such as increased heart rate, loss of appetite and sleep, and an intense feeling of excitement. Culture Humanity's unprecedented set of intellectual skills were a key factor in the species' eventual technological advancement and concomitant domination of the biosphere. Disregarding extinct hominids, humans are the only animals known to teach generalizable information, innately deploy recursive embedding to generate and communicate complex concepts, engage in the "folk physics" required for competent tool design, or cook food in the wild. Teaching and learning preserves the cultural and ethnographic identity of human societies. Other traits and behaviors that are mostly unique to humans include starting fires, phoneme structuring and vocal learning. Language While many species communicate, language is unique to humans, a defining feature of humanity, and a cultural universal. Unlike the limited systems of other animals, human language is openan infinite number of meanings can be produced by combining a limited number of symbols. Human language also has the capacity of displacement, using words to represent things and happenings that are not presently or locally occurring but reside in the shared imagination of interlocutors. Language differs from other forms of communication in that it is modality independent; the same meanings can be conveyed through different media, audibly in speech, visually by sign language or writing, and through tactile media such as braille. Language is central to the communication between humans, and to the sense of identity that unites nations, cultures and ethnic groups. There are approximately six thousand different languages currently in use, including sign languages, and many thousands more that are extinct. The arts Human arts can take many forms including visual, literary, and performing. Visual art can range from paintings and sculptures to film, fashion design, and architecture. Literary arts can include prose, poetry, and dramas. The performing arts generally involve theatre, music, and dance. Humans often combine the different forms (for example, music videos). Other entities that have been described as having artistic qualities include food preparation, video games, and medicine. As well as providing entertainment and transferring knowledge, the arts are also used for political purposes. Art is a defining characteristic of humans and there is evidence for a relationship between creativity and language. The earliest evidence of art was shell engravings made by Homo erectus 300,000 years before modern humans evolved. Art attributed to H. sapiens existed at least 75,000 years ago, with jewellery and drawings found in caves in South Africa. There are various hypotheses as to why humans have adapted to the arts. These include allowing them to better problem solve issues, providing a means to control or influence other humans, encouraging cooperation and contribution within a society or increasing the chance of attracting a potential mate. The use of imagination developed through art, combined with logic may have given early humans an evolutionary advantage. Evidence of humans engaging in musical activities predates cave art and so far music has been practiced by virtually all known human cultures. There exists a wide variety of music genres and ethnic musics; with humans' musical abilities being related to other abilities, including complex social human behaviours. It has been shown that human brains respond to music by becoming synchronized with the rhythm and beat, a process called entrainment. Dance is also a form of human expression found in all cultures and may have evolved as a way to help early humans communicate. Listening to music and observing dance stimulates the orbitofrontal cortex and other pleasure sensing areas of the brain. Unlike speaking, reading and writing does not come naturally to humans and must be taught. Still, literature has been present before the invention of words and language, with 30,000-year-old paintings on walls inside some caves portraying a series of dramatic scenes. One of the oldest surviving works of literature is the Epic of Gilgamesh, first engraved on ancient Babylonian tablets about 4,000 years ago. Beyond simply passing down knowledge, the use and sharing of imaginative fiction through stories might have helped develop humans' capabilities for communication and increased the likelihood of securing a mate. Storytelling may also be used as a way to provide the audience with moral lessons and encourage cooperation. Tools and technologies Stone tools were used by proto-humans at least 2.5 million years ago. The use and manufacture of tools has been put forward as the ability that defines humans more than anything else and has historically been seen as an important evolutionary step. The technology became much more sophisticated about 1.8 million years ago, with the controlled use of fire beginning around 1 million years ago. The wheel and wheeled vehicles appeared simultaneously in several regions some time in the fourth millennium BC. The development of more complex tools and technologies allowed land to be cultivated and animals to be domesticated, thus proving essential in the development of agriculturewhat is known as the Neolithic Revolution. China developed paper, the printing press, gunpowder, the compass and other important inventions. The continued improvements in smelting allowed forging of copper, bronze, iron and eventually steel, which is used in railways, skyscrapers and many other products. This coincided with the Industrial Revolution, where the invention of automated machines brought major changes to humans' lifestyles. Modern technology is observed as progressing exponentially, with major innovations in the 20th century including: electricity, penicillin, semiconductors, internal combustion engines, the Internet, nitrogen fixing fertilizers, airplanes, computers, automobiles, contraceptive pills, nuclear fission, the green revolution, radio, scientific plant breeding, rockets, air conditioning, television and the assembly line. Religion and spirituality Definitions of religion vary; according to one definition, a religion is a belief system concerning the supernatural, sacred or divine, and practices, values, institutions and rituals associated with such belief. Some religions also have a moral code. The evolution and the history of the first religions have become areas of active scientific investigation. Credible evidence of religious behaviour dates to the Middle Paleolithic era (45–200 thousand years ago). It may have evolved to play a role in helping enforce and encourage cooperation between humans. Religion manifests in diverse forms. Religion can include a belief in life after death, the origin of life, the nature of the universe (religious cosmology) and its ultimate fate (eschatology), and moral or ethical teachings. Views on transcendence and immanence vary substantially; traditions variously espouse monism, deism, pantheism, and theism (including polytheism and monotheism). Although measuring religiosity is difficult, a majority of humans profess some variety of religious or spiritual belief. In 2015 the plurality were Christian followed by Muslims, Hindus and Buddhists. As of 2015, about 16%, or slightly under 1.2 billion humans, were irreligious, including those with no religious beliefs or no identity with any religion. Science and philosophy An aspect unique to humans is their ability to transmit knowledge from one generation to the next and to continually build on this information to develop tools, scientific laws and other advances to pass on further. This accumulated knowledge can be tested to answer questions or make predictions about how the universe functions and has been very successful in advancing human ascendancy. Aristotle has been described as the first scientist, and preceded the rise of scientific thought through the Hellenistic period. Other early advances in science came from the Han dynasty in China and during the Islamic Golden Age. The scientific revolution, near the end of the Renaissance, led to the emergence of modern science. A chain of events and influences led to the development of the scientific method, a process of observation and experimentation that is used to differentiate science from pseudoscience. An understanding of mathematics is unique to humans, although other species of animals have some numerical cognition. All of science can be divided into three major branches, the formal sciences (e.g., logic and mathematics), which are concerned with formal systems, the applied sciences (e.g., engineering, medicine), which are focused on practical applications, and the empirical sciences, which are based on empirical observation and are in turn divided into natural sciences (e.g., physics, chemistry, biology) and social sciences (e.g., psychology, economics, sociology). Philosophy is a field of study where humans seek to understand fundamental truths about themselves and the world in which they live. Philosophical inquiry has been a major feature in the development of humans' intellectual history. It has been described as the "no man's land" between definitive scientific knowledge and dogmatic religious teachings. Major fields of philosophy include metaphysics, epistemology, logic, and axiology (which includes ethics and aesthetics). Society Society is the system of organizations and institutions arising from interaction between humans. Humans are highly social and tend to live in large complex social groups. They can be divided into different groups according to their income, wealth, power, reputation and other factors. The structure of social stratification and the degree of social mobility differs, especially between modern and traditional societies. Human groups range from the size of families to nations. The first form of human social organization is thought to have resembled hunter-gatherer band societies. Gender Human societies typically exhibit gender identities and gender roles that distinguish between masculine and feminine characteristics and prescribe the range of acceptable behaviours and attitudes for their members based on their sex. The most common categorisation is a gender binary of men and women. Some societies recognize a third gender, or less commonly a fourth or fifth. In some other societies, non-binary is used as an umbrella term for a range of gender identities that are not solely male or female. Gender roles are often associated with a division of norms, practices, dress, behavior, rights, duties, privileges, status, and power, with men enjoying more rights and privileges than women in most societies, both today and in the past. As a social construct, gender roles are not fixed and vary historically within a society. Challenges to predominant gender norms have recurred in many societies. Little is known about gender roles in the earliest human societies. Early modern humans probably had a range of gender roles similar to that of modern cultures from at least the Upper Paleolithic, while the Neanderthals were less sexually dimorphic and there is evidence that the behavioural difference between males and females was minimal. Kinship All human societies organize, recognize and classify types of social relationships based on relations between parents, children and other descendants (consanguinity), and relations through marriage (affinity). There is also a third type applied to godparents or adoptive children (fictive). These culturally defined relationships are referred to as kinship. In many societies, it is one of the most important social organizing principles and plays a role in transmitting status and inheritance. All societies have rules of incest taboo, according to which marriage between certain kinds of kin relations is prohibited, and some also have rules of preferential marriage with certain kin relations. Pair bonding is a ubiquitous feature of human sexual relationships, whether it is manifested as serial monogamy, polygyny, or polyandry. Genetic evidence indicates that humans were predominantly polygynous for most of their existence as a species, but that this began to shift during the Neolithic, when monogamy started becoming widespread concomitantly with the transition from nomadic to sedentary societies. Anatomical evidence in the form of second-to-fourth digit ratios, a biomarker for prenatal androgen effects, likewise indicates modern humans were polygynous during the Pleistocene. Ethnicity Human ethnic groups are a social category that identifies together as a group based on shared attributes that distinguish them from other groups. These can be a common set of traditions, ancestry, language, history, society, culture, nation, religion, or social treatment within their residing area. Ethnicity is separate from the concept of race, which is based on physical characteristics, although both are socially constructed. Assigning ethnicity to a certain population is complicated, as even within common ethnic designations there can be a diverse range of subgroups, and the makeup of these ethnic groups can change over time at both the collective and individual level. Also, there is no generally accepted definition of what constitutes an ethnic group. Ethnic groupings can play a powerful role in the social identity and solidarity of ethnopolitical units. This has been closely tied to the rise of the nation state as the predominant form of political organization in the 19th and 20th centuries. Government and politics As farming populations gathered in larger and denser communities, interactions between these different groups increased. This led to the development of governance within and between the communities. Humans have evolved the ability to change affiliation with various social groups relatively easily, including previously strong political alliances, if doing so is seen as providing personal advantages. This cognitive flexibility allows individual humans to change their political ideologies, with those with higher flexibility less likely to support authoritarian and nationalistic stances. Governments create laws and policies that affect the citizens that they govern. There have been many forms of government throughout human history, each having various means of obtaining power and the ability to exert diverse controls on the population. Approximately 47% of humans live in some form of a democracy, 17% in a hybrid regime, and 37% in an authoritarian regime. Many countries belong to international organizations and alliances; the largest of these is the United Nations, with 193 member states. Trade and economics Trade, the voluntary exchange of goods and services, is seen as a characteristic that differentiates humans from other animals and has been cited as a practice that gave Homo sapiens a major advantage over other hominids. Evidence suggests early H. sapiens made use of long-distance trade routes to exchange goods and ideas, leading to cultural explosions and providing additional food sources when hunting was sparse, while such trade networks did not exist for the now extinct Neanderthals. Early trade likely involved materials for creating tools like obsidian. The first truly international trade routes were around the spice trade through the Roman and medieval periods. Early human economies were more likely to be based around gift giving instead of a bartering system. Early money consisted of commodities; the oldest being in the form of cattle and the most widely used being cowrie shells. Money has since evolved into governmental issued coins, paper and electronic money. Human study of economics is a social science that looks at how societies distribute scarce resources among different people. There are massive inequalities in the division of wealth among humans; the eight richest humans are worth the same monetary value as the poorest half of all the human population. Conflict Humans commit violence on other humans at a rate comparable to other primates, but have an increased preference for killing adults, infanticide being more common among other primates. Phylogenetic analysis predicts that 2% of early H. sapiens would be murdered, rising to 12% during the medieval period, before dropping to below 2% in modern times. There is great variation in violence between human populations, with rates of homicide about 0.01% in societies that have legal systems and strong cultural attitudes against violence. The willingness of humans to kill other members of their species en masse through organized conflict (i.e., war) has long been the subject of debate. One school of thought holds that war evolved as a means to eliminate competitors, and has always been an innate human characteristic. Another suggests that war is a relatively recent phenomenon and has appeared due to changing social conditions. While not settled, current evidence indicates warlike predispositions only became common about 10,000 years ago, and in many places much more recently than that. War has had a high cost on human life; it is estimated that during the 20th century, between 167 million and 188 million people died as a result of war. War casualty data is less reliable for pre-medieval times, especially global figures. But compared with any period over the past 600 years, the last ~80 years (post 1946), has seen a very significant drop in global military and civilian death rates due to armed conflict.	Biology and health sciences	null	null
682629	https://en.wikipedia.org/wiki/Element%20%28mathematics%29	Element (mathematics)	In mathematics, an element (or member) of a set is any one of the distinct objects that belong to that set. For example, given a set called containing the first four positive integers (), one could say that "3 is an element of ", expressed notationally as . Sets Writing means that the elements of the set are the numbers 1, 2, 3 and 4. Sets of elements of , for example , are subsets of . Sets can themselves be elements. For example, consider the set . The elements of are not 1, 2, 3, and 4. Rather, there are only three elements of , namely the numbers 1 and 2, and the set . The elements of a set can be anything. For example the elements of the set are the color red, the number 12, and the set . In logical terms, . Notation and terminology The binary relation "is an element of", also called set membership, is denoted by the symbol "∈". Writing means that "x is an element of A". Equivalent expressions are "x is a member of A", "x belongs to A", "x is in A" and "x lies in A". The expressions "A includes x" and "A contains x" are also used to mean set membership, although some authors use them to mean instead "x is a subset of A". Logician George Boolos strongly urged that "contains" be used for membership only, and "includes" for the subset relation only. For the relation ∈ , the converse relation ∈T may be written meaning "A contains or includes x". The negation of set membership is denoted by the symbol "∉". Writing means that "x is not an element of A". The symbol ∈ was first used by Giuseppe Peano, in his 1889 work . Here he wrote on page X: which means The symbol ∈ means is. So is read as a is a certain b; … The symbol itself is a stylized lowercase Greek letter epsilon ("ϵ"), the first letter of the word , which means "is". Examples Using the sets defined above, namely A = {1, 2, 3, 4}, B = {1, 2, {3, 4}} and C = {red, green, blue}, the following statements are true: Cardinality of sets The number of elements in a particular set is a property known as cardinality; informally, this is the size of a set. In the above examples, the cardinality of the set A is 4, while the cardinality of set B and set C are both 3. An infinite set is a set with an infinite number of elements, while a finite set is a set with a finite number of elements. The above examples are examples of finite sets. An example of an infinite set is the set of positive integers . Formal relation As a relation, set membership must have a domain and a range. Conventionally the domain is called the universe denoted U. The range is the set of subsets of U called the power set of U and denoted P(U). Thus the relation is a subset of . The converse relation is a subset of .	Mathematics	Set theory	null
83406	https://en.wikipedia.org/wiki/Paracetamol	Paracetamol	Paracetamol, or acetaminophen, is a non-opioid analgesic and antipyretic agent used to treat fever and mild to moderate pain. It is a widely available over-the-counter drug sold under various brand names, including Tylenol and Panadol. Paracetamol relieves pain in both acute mild migraine and episodic tension headache. At a standard dose, paracetamol slightly reduces fever; it is inferior to ibuprofen in that respect, and the benefits of its use for fever are unclear, particularly in the context of fever of viral origins. The aspirin/paracetamol/caffeine combination also helps with both conditions where the pain is mild and is recommended as a first-line treatment for them. Paracetamol is effective for post-surgical pain, but it is inferior to ibuprofen. The paracetamol/ibuprofen combination provides further increase in potency and is superior to either drug alone. The pain relief paracetamol provides in osteoarthritis is small and clinically insignificant. The evidence in its favor for the use in low back pain, cancer pain, and neuropathic pain is insufficient. In the short term, paracetamol is safe and effective when used as directed. Short term adverse effects are uncommon and similar to ibuprofen, but paracetamol is typically safer than nonsteroidal anti-inflammatory drugs (NSAIDs) for long-term use. Paracetamol is also often used in patients who cannot tolerate NSAIDs like ibuprofen. Chronic consumption of paracetamol may result in a drop in hemoglobin level, indicating possible gastrointestinal bleeding, and abnormal liver function tests. The recommended maximum daily dose for an adult is three to four grams. Higher doses may lead to toxicity, including liver failure. Paracetamol poisoning is the foremost cause of acute liver failure in the Western world, and accounts for most drug overdoses in the United States, the United Kingdom, Australia, and New Zealand. Paracetamol was first made in 1878 by Harmon Northrop Morse or possibly in 1852 by Charles Frédéric Gerhardt. It is the most commonly used medication for pain and fever in both the United States and Europe. It is on the World Health Organization's List of Essential Medicines. Paracetamol is available as a generic medication, with brand names including Tylenol and Panadol among others. In 2022, it was the 114th most commonly prescribed medication in the United States, with more than 5million prescriptions. Medical uses Fever Paracetamol is used for reducing fever. However, there has been a lack of research on its antipyretic properties, particularly in adults, and thus its benefits are unclear. As a result, it has been described as over-prescribed for this application. In addition, low-quality clinical data indicates that when used for the common cold, paracetamol may relieve a stuffed or runny nose, but not other cold symptoms such as sore throat, malaise, sneezing, or cough. For people in critical care, paracetamol decreases body temperature by only 0.20.3°C more than control interventions and has no effect on their mortality. It did not change the outcome in febrile patients with stroke. The results are contradictory for paracetamol use in sepsis: higher mortality, lower mortality, and no change in mortality were all reported. Paracetamol offered no benefit in the treatment of dengue fever and was accompanied by a higher rate of liver enzyme elevation: a sign of potential liver damage. Overall, there is no support for a routine administration of antipyretic drugs, including paracetamol, to hospitalized patients with fever and infection. The efficacy of paracetamol in children with fever is unclear. Paracetamol should not be used solely to reduce body temperature; however, it may be considered for children with fever who appear distressed. It does not prevent febrile seizures. It appears that 0.2°C decrease of the body temperature in children after a standard dose of paracetamol is of questionable value, particularly in emergencies. Based on this, some physicians advocate using higher doses that may decrease the temperature by as much as 0.7°C. Meta-analyses showed that paracetamol is less effective than ibuprofen in children (marginally less effective, according to another analysis), including children younger than 2 years old, with equivalent safety. Exacerbation of asthma occurs with similar frequency for both medications. Giving paracetamol and ibuprofen together at the same time to children under 5 is not recommended; however, doses may be alternated if required. Pain Paracetamol is used for the relief of mild to moderate pain such as headache, muscle aches, minor arthritis pain, toothache as well as pain caused by cold, flu, sprains, and dysmenorrhea. It is recommended, in particular, for acute mild to moderate pain, since the evidence for the treatment of chronic pain is insufficient. Musculoskeletal pain The benefits of paracetamol in musculoskeletal conditions, such as osteoarthritis and backache, are uncertain. It appears to provide only small and not clinically important benefits in osteoarthritis. American College of Rheumatology and Arthritis Foundation guideline for the management of osteoarthritis notes that the effect size in clinical trials of paracetamol has been very small, which suggests that for most individuals it is ineffective. The guideline conditionally recommends paracetamol for short-term and episodic use to those who do not tolerate nonsteroidal anti-inflammatory drugs. For people taking it regularly, monitoring for liver toxicity is required. Essentially the same recommendation was issued by EULAR for hand osteoarthritis. Similarly, the ESCEO algorithm for the treatment of knee osteoarthritis recommends limiting the use of paracetamol to short-term rescue analgesia only. Paracetamol is ineffective for acute low back pain. No randomized clinical trials evaluated its use for chronic or radicular back pain, and the evidence in favor of paracetamol is lacking. Headaches Paracetamol is effective for acute migraine: 39 % of people experience pain relief at one hour compared with 20 % in the control group. The aspirin/paracetamol/caffeine combination also "has strong evidence of effectiveness and can be used as a first-line treatment for migraine". Paracetamol on its own only slightly alleviates episodic tension headache in those who have them frequently. However, the aspirin/paracetamol/caffeine combination is superior to both paracetamol alone and placebo and offers meaningful relief of tension headache: two hours after administering the medication, 29 % of those who took the combination were pain-free as compared with 21 % on paracetamol and 18 % on placebo. The German, Austrian, and Swiss headache societies and the German Society of Neurology recommend this combination as a "highlighted" one for self-medication of tension headache, with paracetamol/caffeine combination being a "remedy of first choice", and paracetamol a "remedy of second choice". Dental and other post-surgical pain Pain after a dental surgery provides a reliable model for the action of analgesics on other kinds of acute pain. For the relief of such pain, paracetamol is inferior to ibuprofen. Full therapeutic doses of nonsteroidal anti-inflammatory drugs (NSAIDs) ibuprofen, naproxen or diclofenac are clearly more efficacious than the paracetamol/codeine combination which is frequently prescribed for dental pain. The combinations of paracetamol and NSAIDs ibuprofen or diclofenac are promising, possibly offering better pain control than either paracetamol or the NSAID alone. Additionally, the paracetamol/ibuprofen combination may be superior to paracetamol/codeine and ibuprofen/codeine combinations. A meta-analysis of general post-surgical pain, which included dental and other surgery, showed the paracetamol/codeine combination to be more effective than paracetamol alone: it provided significant pain relief to as much as 53 % of the participants, while the placebo helped only 7 %. Other pain Paracetamol fails to relieve procedural pain in newborn babies. For perineal pain postpartum paracetamol appears to be less effective than nonsteroidal anti-inflammatory drugs (NSAIDs). The studies to support or refute the use of paracetamol for cancer pain and neuropathic pain are lacking. There is limited evidence in favor of the use of the intravenous form of paracetamol for acute pain control in the emergency department. The combination of paracetamol with caffeine is superior to paracetamol alone for the treatment of acute pain. Patent ductus arteriosus Paracetamol helps ductal closure in patent ductus arteriosus. It is as effective for this purpose as ibuprofen or indomethacin, but results in less frequent gastrointestinal bleeding than ibuprofen. Its use for extremely low birth weight and gestational age infants however requires further study. Adverse effects Gastrointestinal adverse effects such as nausea and abdominal pain are extremely uncommon, and their frequency is nothing like that of ibuprofen. Increase in risk-taking behavior is possible. According to the U.S. Food and Drug Administration (FDA), the drug may cause rare and possibly fatal skin reactions such as Stevens–Johnson syndrome and toxic epidermal necrolysis, Rechallenge tests and an analysis of American but not French pharmacovigilance databases indicated a risk of these reactions. In clinical trials for osteoarthritis, the number of participants reporting adverse effects was similar for those on paracetamol and on placebo. However, the abnormal liver function tests (meaning there was some inflammation or damage to the liver) were almost four times more likely in those on paracetamol, although the clinical importance of this effect is uncertain. After 13 weeks of paracetamol therapy for knee pain, a drop in hemoglobin level indicating gastrointestinal bleeding was observed in 20 % of participants, this rate being similar to the ibuprofen group. Due to the absence of controlled studies, most of the information about the long-term safety of paracetamol comes from observational studies. These indicate a consistent pattern of increased mortality as well as cardiovascular (stroke, myocardial infarction), gastrointestinal (ulcers, bleeding) and renal adverse effects with increased dose of paracetamol. Use of paracetamol is associated with 1.9 times higher risk of peptic ulcer. Those who take it regularly at a higher dose (more than 23g daily) are at much higher risk (3.63.7 times) of gastrointestinal bleeding and other bleeding events. Meta-analyses suggest that paracetamol may increase the risk of kidney impairment by 23 % and kidney cancer by 28 %. Paracetamol slightly but significantly increases blood pressure and heart rate. A 2022 double-blind, placebo-controlled, crossover study has provided evidence that daily, high-dose use (4 g per day) of paracetamol increases systolic BP. A review of available research has suggested that increase in systolic blood pressure and increased risk of gastrointestinal bleeding associated with chronic paracetamol use shows a degree of dose dependence. The association between paracetamol use and asthma in children has been a matter of controversy. However, the most recent research suggests that there is no association, and that the frequency of asthma exacerbations in children after paracetamol is the same as after another frequently used pain killer, ibuprofen. In recommended doses, the side effects of paracetamol are mild to non-existent. In contrast to aspirin, it is not a blood thinner (and thus may be used in patients where bleeding is a concern), and it does not cause gastric irritation. Compared to Ibuprofen—which can have adverse effects that include diarrhea, vomiting, and abdominal pain—paracetamol is well tolerated with fewer side effects. Prolonged daily use may cause kidney or liver damage. Paracetamol is metabolized by the liver and is hepatotoxic; side effects may be more likely in chronic alcoholics or patients with liver damage. Until 2010 paracetamol was believed safe in pregnancy however, in a study published in October 2010 it has been linked to infertility in the adult life of the unborn. Like NSAIDs and unlike opioid analgesics, paracetamol has not been found to cause euphoria or alter mood. One recent research study showed evidence that paracetamol can ease psychological pain, but more studies are needed to draw a stronger conclusion. Unlike aspirin, it is safe for children, as paracetamol is not associated with a risk of Reye's syndrome in children with viral illnesses. Chronic users of paracetamol may have a higher risk of developing blood cancer. Use in pregnancy Paracetamol safety in pregnancy has been under increased scrutiny. There appears to be no link between paracetamol use in the first trimester and adverse pregnancy outcomes or birth defects. However, indications exist of a possible increase in the risk of asthma and developmental and reproductive disorders in the offspring of women with prolonged use of paracetamol during pregnancy. Paracetamol use by the mother during pregnancy is associated with an increased risk of childhood asthma, but so are the maternal infections for which paracetamol may be used, and separating these influences is difficult. Paracetamol, in a small-scale meta-analysis was also associated with a 2030 % increase in autism spectrum disorder, attention deficit hyperactivity disorder, and conduct disorder, with the association being lower in a meta-analysis where a larger demographic was used, but it is unclear whether this is a causal relationship and whether there was potential bias in the findings. There is also an argument that the large number, consistency, and robust designs of the studies provide strong evidence in favor of paracetamol causing the increased risk of these neurodevelopmental disorders. In animal experiments, paracetamol disrupts fetal testosterone production, and several epidemiological studies linked cryptorchidism with mother's paracetamol use for more than two weeks in the second trimester. On the other hand, several studies did not find any association. The consensus recommendation appears to be to avoid prolonged use of paracetamol in pregnancy and use it only when necessary, at the lowest effective dosage, and for the shortest time. In pregnancy, paracetamol and metoclopramide are deemed safe as are NSAIDs until the third trimester. Overdose Overdose of paracetamol is caused by taking more than the recommended maximum daily dose of paracetamol for healthy adults (three or four grams), and can cause potentially fatal liver damage. A single dose should not exceed 1000 mg, doses should be taken no sooner than four hours apart, and no more than four doses (4000 mg) in 24 hours. While a majority of adult overdoses are linked to suicide attempts, many cases are accidental, often due to the use of more than one paracetamol-containing product over an extended period. Paracetamol toxicity has become the foremost cause of acute liver failure in the United States by 2003, and , paracetamol accounted for most drug overdoses in the United States, the United Kingdom, Australia, and New Zealand. As of 2004, paracetamol overdose resulted in more calls to poison control centers in the U.S. than overdose of any other pharmacological substance. According to the FDA, in the United States, "56,000 emergency room visits, 26,000 hospitalizations, and 458 deaths per year [were] related to acetaminophen-associated overdoses during the 1990s. Within these estimates, unintentional acetaminophen overdose accounted for nearly 25 % of the emergency department visits, 10 % of the hospitalizations, and 25 % of the deaths." Overdoses are frequently related to high-dose recreational use of prescription opioids, as these opioids are most often combined with paracetamol. The overdose risk may be heightened by frequent consumption of alcohol. Untreated paracetamol overdose results in a lengthy, painful illness. Signs and symptoms of paracetamol toxicity may initially be absent or non-specific symptoms. The first symptoms of overdose usually begin several hours after ingestion, with nausea, vomiting, sweating, and pain as acute liver failure starts. People who take overdoses of paracetamol do not fall asleep or lose consciousness, although most people who attempt suicide with paracetamol wrongly believe that they will be rendered unconscious by the drug. Treatment is aimed at removing the paracetamol from the body and replenishing glutathione. Activated charcoal can be used to decrease absorption of paracetamol if the person comes to the hospital soon after the overdose. While the antidote, acetylcysteine (also called N-acetylcysteine or NAC), acts as a precursor for glutathione, helping the body regenerate enough to prevent or at least decrease the possible damage to the liver; a liver transplant is often required if damage to the liver becomes severe. NAC was usually given following a treatment nomogram (one for people with risk factors, and one for those without), but the use of the nomogram is no longer recommended as evidence to support the use of risk factors was poor and inconsistent, and many of the risk factors are imprecise and difficult to determine with sufficient certainty in clinical practice. Toxicity of paracetamol is due to its quinone metabolite NAPQI and NAC also helps in neutralizing it. Kidney failure is also a possible side effect. Interactions Prokinetic agents such as metoclopramide accelerate gastric emptying, shorten time (tmax) to paracetamol peak blood plasma concentration (Cmax), and increase Cmax. Medications slowing gastric emptying such as propantheline and morphine lengthen tmax and decrease Cmax. The interaction with morphine may result in patients failing to achieve the therapeutic concentration of paracetamol; the clinical significance of interactions with metoclopramide and propantheline is unclear. There have been suspicions that cytochrome inducers may enhance the toxic pathway of paracetamol metabolism to NAPQI (see Paracetamol#Pharmacokinetics). By and large, these suspicions have not been confirmed. Out of the inducers studied, the evidence of potentially increased liver toxicity in paracetamol overdose exists for phenobarbital, primidone, isoniazid, and possibly St John's wort. On the other hand, the anti-tuberculosis drug isoniazid cuts the formation of NAPQI by 70%. Ranitidine increased paracetamol area under the curve (AUC) 1.6-fold. AUC increases are also observed with nizatidine and cisapride. The effect is explained by these drugs inhibiting glucuronidation of paracetamol. Paracetamol raises plasma concentrations of ethinylestradiol by 22 % by inhibiting its sulfation. Paracetamol increases INR during warfarin therapy and should be limited to no more than 2 g per week. Pharmacology Pharmacodynamics Paracetamol appears to exert its effects through two mechanisms: the inhibition of cyclooxygenase (COX) and actions of its metabolite N-arachidonoylphenolamine (AM404). Supporting the first mechanism, pharmacologically and in its side effects, paracetamol is close to classical nonsteroidal anti-inflammatory drugs (NSAIDs) that act by inhibiting COX-1 and COX-2 enzymes and especially similar to selective COX-2 inhibitors. Paracetamol inhibits prostaglandin synthesis by reducing the active form of COX-1 and COX-2 enzymes. This occurs only when the concentration of arachidonic acid and peroxides is low. Under these conditions, COX-2 is the predominant form of cyclooxygenase, which explains the apparent COX-2 selectivity of paracetamol. Under the conditions of inflammation, the concentration of peroxides is high, which counteracts the reducing effect of paracetamol. Accordingly, the anti-inflammatory action of paracetamol is slight. The anti-inflammatory action of paracetamol (via COX inhibition) has also been found to primarily target the central nervous system and not peripheral areas of the body, explaining the lack of side effects associated with conventional NSAIDs such as gastric bleeding. The second mechanism centers on the paracetamol metabolite AM404. This metabolite has been detected in the brains of animals and cerebrospinal fluid of humans taking paracetamol. It is formed in the brain from another paracetamol metabolite 4-aminophenol by action of fatty acid amide hydrolase. AM404 is a weak agonist of cannabinoid receptors CB1 and CB2, an inhibitor of endocannabinoid transporter, and a potent activator of TRPV1 receptor. This and other research indicate that the endocannabinoid system and TRPV1 may play an important role in the analgesic effect of paracetamol. In 2018, Suemaru et al. found that, in mice, paracetamol exerts an anticonvulsant effect by activation of the TRPV1 receptors and a decrease in neuronal excitability by hyperpolarization of neurons. The exact mechanism of the anticonvulsant effect of acetaminophen is not clear. According to Suemaru et al., acetaminophen and its active metabolite AM404 show a dose-dependent anticonvulsant activity against pentylenetetrazol-induced seizures in mice. Pharmacokinetics After being taken by mouth, paracetamol is rapidly absorbed from the small intestine, while absorption from the stomach is negligible. Thus, the rate of absorption depends on stomach emptying. Food slows the stomach's emptying and absorption, but the total amount absorbed stays the same. In the same subjects, the peak plasma concentration of paracetamol was reached after 20 minutes when fasting versus 90 minutes when fed. High carbohydrate (but not high protein or high fat) food decreases paracetamol peak plasma concentration by four times. Even in the fasting state, the rate of absorption of paracetamol is variable and depends on the formulation, with maximum plasma concentration being reached after 20 minutes to 1.5 hours. Paracetamol's bioavailability is dose-dependent: it increases from 63 % for 500mg dose to 89 % for 1000mg dose. Its plasma terminal elimination half-life is 1.92.5 hours, and volume of distribution is roughly 50L. Protein binding is negligible, except under the conditions of overdose, when it may reach 1521 %. The concentration in serum after a typical dose of paracetamol usually peaks below 30μg/mL (200μmol/L). After 4 hours, the concentration is usually less than 10μg/mL (66μmol/L). Paracetamol is metabolized primarily in the liver, mainly by glucuronidation and sulfation, and the products are then eliminated in the urine (see the Scheme on the right). Only 25 % of the drug is excreted unchanged in the urine. Glucuronidation by UGT1A1 and UGT1A6 accounts for 5070 % of the drug metabolism. Additional 2535 % of paracetamol is converted to sulfate by sulfation enzymes SULT1A1, SULT1A3, and SULT1E1. A minor metabolic pathway (5–15 %) of oxidation by cytochrome P450 enzymes, mainly by CYP2E1, forms a toxic metabolite known as NAPQI (N-acetyl-p-benzoquinone imine). NAPQI is responsible for the liver toxicity of paracetamol. At usual doses of paracetamol, NAPQI is quickly detoxified by conjugation with glutathione. The non-toxic conjugate APAP-GSH is taken up in the bile and further degraded to mercapturic and cysteine conjugates that are excreted in the urine. In overdose, glutathione is depleted by a large amount of formed NAPQI, and NAPQI binds to mitochondria proteins of the liver cells causing oxidative stress and toxicity. Yet another minor but important direction of metabolism is deacetylation of 12 % of paracetamol to form p-aminophenol. p-Aminophenol is then converted in the brain by fatty acid amide hydrolase into AM404, a compound that may be partially responsible for the analgesic action of paracetamol. Chemistry Synthesis Classical methods The classical methods for the production of paracetamol involve the acetylation of 4-aminophenol with acetic anhydride as the last step. They differ in how 4-aminophenol is prepared. In one method, nitration of phenol with nitric acid affords 4-nitrophenol, which is reduced to 4-aminophenol by hydrogenation over Raney nickel. In another method, nitrobenzene is reduced electrolytically giving 4-aminophenol directly. Additionally, 4-nitrophenol can be selectively reduced by Tin(II) Chloride in absolute ethanol or ethyl acetate to produce a 91 % yield of 4-aminophenol. Celanese synthesis An alternative industrial synthesis developed at Celanese involves firstly direct acylation of phenol with acetic anhydride in the presence of hydrogen fluoride to a ketone, then the conversion of the ketone with hydroxylamine to a ketoxime, and finally the acid-catalyzed Beckmann rearrangement of the cetoxime to the para-acetylaminophenol product. Reactions 4-Aminophenol may be obtained by the amide hydrolysis of paracetamol. This reaction is also used to determine paracetamol in urine samples: After hydrolysis with hydrochloric acid, 4-aminophenol reacts in ammonia solution with a phenol derivate, e.g. salicylic acid, to form an indophenol dye under oxidization by air. History Acetanilide was the first aniline derivative serendipitously found to possess analgesic as well as antipyretic properties, and was quickly introduced into medical practice under the name of Antifebrin by Cahn & Hepp in 1886. But its unacceptable toxic effectsthe most alarming being cyanosis due to methemoglobinemia, an increase of hemoglobin in its ferric [Fe3+] state, called methemoglobin, which cannot bind oxygen, and thus decreases overall carriage of oxygen to tissueprompted the search for less toxic aniline derivatives. Some reports state that Cahn & Hepp or a French chemist called Charles Gerhardt first synthesized paracetamol in 1852. Harmon Northrop Morse synthesized paracetamol at Johns Hopkins University via the reduction of p-nitrophenol with tin in glacial acetic acid in 1877, but it was not until 1887 that clinical pharmacologist Joseph von Mering tried paracetamol on humans. In 1893, von Mering published a paper reporting on the clinical results of paracetamol with phenacetin, another aniline derivative. Von Mering claimed that, unlike phenacetin, paracetamol had a slight tendency to produce methemoglobinemia. Paracetamol was then quickly discarded in favor of phenacetin. The sales of phenacetin established Bayer as a leading pharmaceutical company. Von Mering's claims remained essentially unchallenged for half a century until two teams of researchers from the United States analyzed the metabolism of acetanilide and phenacetin. In 1947, David Lester and Leon Greenberg found strong evidence that paracetamol was a major metabolite of acetanilide in human blood, and in a subsequent study they reported that large doses of paracetamol given to albino rats did not cause methemoglobinemia. In 1948, Bernard Brodie, Julius Axelrod and Frederick Flinn confirmed that paracetamol was the major metabolite of acetanilide in humans, and established that it was just as efficacious an analgesic as its precursor. They also suggested that methemoglobinemia is produced in humans mainly by another metabolite, phenylhydroxylamine. A follow-up paper by Brodie and Axelrod in 1949 established that phenacetin was also metabolized to paracetamol. This led to a "rediscovery" of paracetamol. Paracetamol was first marketed in the United States in 1950 under the name Trigesic, a combination of paracetamol, aspirin, and caffeine. Reports in 1951 of three users stricken with the blood disease agranulocytosis led to its removal from the marketplace, and it took several years until it became clear that the disease was unconnected. The following year, 1952, paracetamol returned to the U.S. market as a prescription drug. In the United Kingdom, marketing of paracetamol began in 1956 by Sterling-Winthrop Co. as Panadol, available only by prescription, and promoted as preferable to aspirin since it was safe for children and people with ulcers. In 1963, paracetamol was added to the British Pharmacopoeia, and has gained popularity since then as an analgesic agent with few side-effects and little interaction with other pharmaceutical agents. Concerns about paracetamol's safety delayed its widespread acceptance until the 1970s, but in the 1980s paracetamol sales exceeded those of aspirin in many countries, including the United Kingdom. This was accompanied by the commercial demise of phenacetin, blamed as the cause of analgesic nephropathy and hematological toxicity. Available in the U.S. without a prescription since 1955 (1960, according to another source), paracetamol has become a common household drug. In 1988, Sterling Winthrop was acquired by Eastman Kodak which sold the over the counter drug rights to SmithKline Beecham in 1994. In June 2009, an FDA advisory committee recommended that new restrictions be placed on paracetamol use in the United States to help protect people from the potential toxic effects. The maximum single adult dosage would be decreased from 1000mg to 650mg, while combinations of paracetamol and other products would be prohibited. Committee members were particularly concerned by the fact that the then-present maximum dosages of paracetamol had been shown to produce alterations in liver function. In January 2011, the FDA asked manufacturers of prescription combination products containing paracetamol to limit its amount to no more than 325mg per tablet or capsule and began requiring manufacturers to update the labels of all prescription combination paracetamol products to warn of the potential risk of severe liver damage. Manufacturers had three years to limit the amount of paracetamol in their prescription drug products to 325mg per dosage unit. In November 2011, the Medicines and Healthcare products Regulatory Agency revised UK dosing of liquid paracetamol for children. In September 2013, "Use Only as Directed", an episode of the radio program This American Life highlighted deaths from paracetamol overdose. This report was followed by two reports by ProPublica alleging that the "FDA has long been aware of studies showing the risks of acetaminophen. So has the maker of Tylenol, McNeil Consumer Healthcare, a division of Johnson & Johnson" and "McNeil, the maker of Tylenol, ... has repeatedly opposed safety warnings, dosage restrictions and other measures meant to safeguard users of the drug." Society and culture Naming Paracetamol is the Australian Approved Name and British Approved Name as well as the international nonproprietary name used by the WHO and in many other countries; acetaminophen is the United States Adopted Name and Japanese Accepted Name and also the name generally used in Canada, Venezuela, Colombia, and Iran. Both paracetamol and acetaminophen are contractions of chemical names for the compound. The word "paracetamol" is a shortened form of para-acetylaminophenol, and was coined by Frederick Stearns & Co in 1956, while the word "acetaminophen" is a shortened form of N-acetyl-p-aminophenol (APAP), which was coined and first marketed by McNeil Laboratories in 1955. The initialism APAP is used by dispensing pharmacists in the United States. Available forms Paracetamol is available in oral, suppository, and intravenous forms. Intravenous paracetamol is sold under the brand name Ofirmev in the United States. In some formulations, paracetamol is combined with the opiate codeine, sometimes referred to as co-codamol (BAN) and Panadeine in Australia. In the U.S., this combination is available only by prescription. As of 1 February 2018, medications containing codeine also became prescription-only in Australia. Paracetamol is also combined with other opioids such as dihydrocodeine, referred to as co-dydramol (British Approved Name (BAN)), oxycodone or hydrocodone. Another very commonly used analgesic combination includes paracetamol in combination with propoxyphene napsylate. A combination of paracetamol, codeine, and the doxylamine succinate is also available. Paracetamol is sometimes combined with phenylephrine hydrochloride. Sometimes a third active ingredient, such as ascorbic acid, caffeine, chlorpheniramine maleate, or guaifenesin is added to this combination. Research Claims that paracetamol is an effective analgesic medication to treat symptoms of COVID-19 were found to be unsubstantiated. Veterinary use Cats Paracetamol is extremely toxic to cats, which lack the necessary UGT1A6 enzyme to detoxify it. Initial symptoms include vomiting, salivation, and discoloration of the tongue and gums. Unlike an overdose in humans, liver damage is rarely the cause of death; instead, methemoglobin formation and the production of Heinz bodies in red blood cells inhibit oxygen transport by the blood, causing asphyxiation (methemoglobinemia and hemolytic anemia). Treatment of the toxicosis with acetylcysteine is recommended. Dogs Paracetamol has been reported to be as effective as aspirin in the treatment of musculoskeletal pain in dogs. A paracetamol–codeine product (brand name Pardale-V) licensed for use in dogs is available for purchase under supervision of a vet, pharmacist or other qualified person. It should be administered to dogs only on veterinary advice and with extreme caution. The main effect of toxicity in dogs is liver damage, and GI ulceration has been reported. Acetylcysteine treatment is efficacious in dogs when administered within two hours of paracetamol ingestion. Snakes Paracetamol is lethal to snakes and has been suggested as a chemical control program for the invasive brown tree snake (Boiga irregularis) in Guam. Doses of 80mg are inserted into dead mice that are scattered by helicopter as lethal bait to be consumed by the snakes.	Biology and health sciences	Pain treatments	Health
83443	https://en.wikipedia.org/wiki/Birth	Birth	Birth is the act or process of bearing or bringing forth offspring, also referred to in technical contexts as parturition. In mammals, the process is initiated by hormones which cause the muscular walls of the uterus to contract, expelling the fetus at a developmental stage when it is ready to feed and breathe. In some species, the offspring is precocial and can move around almost immediately after birth but in others, it is altricial and completely dependent on parenting. In marsupials, the fetus is born at a very immature stage after a short gestation and develops further in its mother's womb pouch. It is not only mammals that give birth. Some reptiles, amphibians, fish and invertebrates carry their developing young inside them. Some of these are ovoviviparous, with the eggs being hatched inside the mother's body, and others are viviparous, with the embryo developing inside their body, as in the case of mammals. Human childbirth Humans usually produce a single offspring at a time. The mother's body is prepared for birth by hormones produced by the pituitary gland, the ovary and the placenta. The total gestation period from fertilization to birth is normally about 38 weeks (birth usually occurring 40 weeks after the last menstrual period). The normal process of childbirth takes several hours and has three stages. The first stage starts with a series of involuntary contractions of the muscular walls of the uterus and gradual dilation of the cervix. The active phase of the first stage starts when the cervix is dilated more than about 4 cm in diameter and is when the contractions become stronger and regular. The head (or the buttocks in a breech birth) of the baby is pushed against the cervix, which gradually dilates until it is fully dilated at 10 cm diameter. At some time, the amniotic sac bursts and the amniotic fluid escapes (also known as rupture of membranes or breaking the water). In stage two, starting when the cervix is fully dilated, strong contractions of the uterus and active pushing by the mother expels the baby out through the vagina, which during this stage of labour is called a birth canal as this passage contains a baby, and the baby is born with umbilical cord attached. In stage three, which begins after the birth of the baby, further contractions expel the placenta, amniotic sac, and the remaining portion of the umbilical cord usually within a few minutes. Enormous changes take place in the newborn's circulation to enable breathing in air. In the uterus, the fetus is dependent on circulation of blood through the placenta for sustenance including gaseous exchange and the unborn baby's blood bypasses the lungs by flowing through the foramen ovale, which is a hole in the septum dividing the right atrium and left atrium. After birth the umbilical cord is clamped and cut, the baby starts to breathe air, and blood from the right ventricle starts to flow to the lungs for gaseous exchange and oxygenated blood returns to the left atrium, which is pumped into the left ventricle, and then pumped into the main arterial system. As a result of these changes, the blood pressure in the left atrium exceeds the pressure in the right atrium, and this pressure difference forces the foramen ovale to close separating the left and right sides of the heart. The umbilical vein, umbilical arteries, ductus venosus and ductus arteriosus are not needed for life in air and in time these vessels become ligaments (embryonic remnants). Other mammals Large mammals, such as primates, cattle, horses, some antelopes, giraffes, hippopotamuses, rhinoceroses, elephants, seals, whales, dolphins, and porpoises, generally are pregnant with one offspring at a time, although they may have twin or multiple births on occasion. In these large animals, the birth process is similar to that of a human, though in most the offspring is precocial. This means that it is born in a more advanced state than a human baby and is able to stand, walk and run (or swim in the case of an aquatic mammal) shortly after birth. In the case of whales, dolphins and porpoises, the single calf is normally born tail first which minimizes the risk of drowning. The mother encourages the newborn calf to rise to the surface of the water to breathe. Large mammals which give birth to twins is much more rare, but it does occur occasionally even for mammals as large as elephants. In April 2018, approximately 8-month old elephant twins were sighted joining their mother's herd in the Tarangire National Park of Tanzania, estimated to have been born in August 2017. Cattle Birthing in cattle is typical of a larger mammal. A cow goes through three stages of labor during normal delivery of a calf. During stage one, the animal seeks a quiet place away from the rest of the herd. Hormone changes cause soft tissues of the birth canal to relax as the mother's body prepares for birth. The contractions of the uterus are not obvious externally, but the cow may be restless. She may appear agitated, alternating between standing and lying down, with her tail slightly raised and her back arched. The fetus is pushed toward the birth canal by each contraction and the cow's cervix gradually begins to dilate. Stage one may last several hours, and ends when the cervix is fully dilated. Stage two can be seen to be underway when there is external protrusion of the amniotic sac through the vulva, closely followed by the appearance of the calf's front hooves and head in a front presentation (or occasionally the calf's tail and rear end in a posterior presentation). During the second stage, the cow will usually lie down on her side to push and the calf progresses through the birth canal. The complete delivery of the calf (or calves in a multiple birth) signifies the end of stage two. The cow scrambles to her feet (if lying down at this stage), turns round and starts vigorously licking the calf. The calf takes its first few breaths and within minutes is struggling to rise to its feet. The third and final stage of labor is the delivery of the placenta, which is usually expelled within a few hours and is often eaten by the normally herbivorous cow. Dogs Birth is termed whelping in dogs. Among dogs, as whelping approaches, contractions become more frequent. Labour in the bitch can be divided into 3 stages. The first stage is when the cervix dilates, causing discomfort and restlessness in the dog. Common signs of this stage are panting, fasting, and/or vomiting. This may last up to 12 hours. Stage two is the passage of the offspring. The amniotic sac looking like a glistening grey balloon, with a puppy inside, is propelled through the vulva. After further contractions, the sac is expelled and the bitch breaks the membranes, releasing clear fluid and exposing the puppy. The mother chews at the umbilical cord and licks the puppy vigorously, which stimulates it to breathe. If the puppy has not taken its first breath within about six minutes, it is likely to die. Further puppies follow in a similar way one by one usually with less straining than the first usually at 15-60-minute intervals. If a pup has not been passed in 2 hours a veterinarian should be contacted. Stage three is the passing of the placentas. This often occurs in conjunction with stage two with the passing of each offspring. The mother will then usually eat the afterbirth. This is an adaption to keep the den clean and prevent its detection by predators. Marsupials An infant marsupial is born in a very immature state. The gestation period is usually shorter than the intervals between oestrus periods. The first sign that a birth is imminent is the mother cleaning out her pouch. When it is born, the infant is pink, blind, furless and a few centimetres long. It has nostrils in order to breathe and forelegs to cling onto its mother's hairs but its hind legs are undeveloped. It crawls through its mother's fur and makes its way into the pouch. Here it fixes onto a teat which swells inside its mouth. It stays attached to the teat for several months until it is sufficiently developed to emerge. Joeys are born with "oral shields"; in species without pouches or with rudimentary pouches these are more developed than in forms with well-developed pouches, implying a role in maintaining the young attached to the mother's nipple. Other animals Many reptiles and the vast majority of invertebrates, most fish, amphibians and all birds are oviparous, that is, they lay eggs with little or no embryonic development taking place within the mother. In aquatic organisms, fertilization is nearly always external with sperm and eggs being liberated into the water (an exception is sharks and rays, which have internal fertilization). Millions of eggs may be produced with no further parental involvement, in the expectation that a small number may survive to become mature individuals. Terrestrial invertebrates may also produce large numbers of eggs, a few of which may avoid predation and carry on the species. Some fish, reptiles, and amphibians have adopted a different strategy and invest their effort in producing a small number of young at a more advanced stage which are more likely to survive to adulthood. Birds care for their young in the nest and provide for their needs after hatching and it is perhaps unsurprising that internal development does not occur in birds, given their need to fly. Ovoviviparity is a mode of reproduction in which embryos develop inside eggs that remain in the mother's body until they are ready to hatch. Ovoviviparous animals are similar to viviparous species in that there is internal fertilization and the young are born in an advanced state, but differ in that there is no placental connection and the unborn young are nourished by egg yolk. The mother's body provides gas exchange (respiration), but that is largely necessary for oviparous animals as well. In many sharks the eggs hatch in the oviduct within the mother's body and the embryos are nourished by the egg's yolk and fluids secreted by glands in the walls of the oviduct. The Lamniforme sharks practice oophagy, where the first embryos to hatch consume the remaining eggs and sand tiger shark pups cannibalistically consume neighbouring embryos. The requiem sharks maintain a placental link to the developing young, this practice is known as viviparity. This is more analogous to mammalian gestation than to that of other fishes. In all these cases, the young are born alive and fully functional. The majority of caecilians are ovoviviparous and give birth to already developed offspring. When the young have finished their yolk sacs they feed on nutrients secreted by cells lining the oviduct and even the cells themselves which they eat with specialist scraping teeth. The Alpine salamander (Salamandra atra) and several species of Tanzanian toad in the genus Nectophrynoides are ovoviviparous, developing through the larval stage inside the mother's oviduct and eventually emerging as fully formed juveniles. A more developed form of viviparity called placental viviparity is adopted by some species of scorpions and cockroaches, certain genera of sharks, snakes and velvet worms. In these, the developing embryo is nourished by some form of placental structure. The earliest known placenta was found recently in a group of extinct fishes called placoderms. A fossil from Australia's Gogo Formation, laid down in the Devonian period, 380 million years ago, was found with an embryo inside it connected by an umbilical cord to a yolk sac. The find confirmed the hypothesis that a sub-group of placoderms, called ptyctodontids, fertilized their eggs internally. Some fishes that fertilize their eggs internally also give birth to live young, as seen here. This discovery moved our knowledge of live birth back 200 million years. The fossil of another genus was found with three embryos in the same position. Placoderms are a sister group of the ancestor of all living jawed fishes (Gnathostomata), including both chondrichthyans, the sharks & rays, and Osteichthyes, the bony fishes. Among lizards, the viviparous lizard Zootoca vivipara, the Jackson's chameleon, slow worms and many species of skink are viviparous, giving birth to live young. Some are ovoviviparous but others such as members of the genera Tiliqua and Corucia, give birth to live young that develop internally, deriving their nourishment from a mammal-like placenta attached to the inside of the mother's uterus. In a recently described example, an African species, Trachylepis ivensi, has developed a purely reptilian placenta directly comparable in structure and function to a mammalian placenta. Vivipary is rare in snakes, but boas and vipers are viviparous, giving birth to live young. The majority of insects lay eggs but a very few give birth to offspring that are miniature versions of the adult. The aphid has a complex life cycle and during the summer months is able to multiply with great rapidity. Its reproduction is typically parthenogenetic and viviparous and females produce unfertilized eggs which they retain within their bodies. The embryos develop within their mothers' ovarioles and the offspring are clones of their mothers. Female nymphs are born which grow rapidly and soon produce more female offspring themselves. In some instances, the newborn nymphs already have developing embryos inside them.	Biology and health sciences	Animal reproduction	null
83537	https://en.wikipedia.org/wiki/Anemia	Anemia	Anemia (also spelled anaemia in British English) is a blood disorder in which the blood has a reduced ability to carry oxygen. This can be due to a lower than normal number of red blood cells, a reduction in the amount of hemoglobin available for oxygen transport, or abnormalities in hemoglobin that impair its function. The name is derived . When anemia comes on slowly, the symptoms are often vague, such as tiredness, weakness, shortness of breath, headaches, and a reduced ability to exercise. When anemia is acute, symptoms may include confusion, feeling like one is going to pass out, loss of consciousness, and increased thirst. Anemia must be significant before a person becomes noticeably pale. Additional symptoms may occur depending on the underlying cause. Anemia can be temporary or long term and can range from mild to severe. Anemia can be caused by blood loss, decreased red blood cell production, and increased red blood cell breakdown. Causes of blood loss include bleeding due to inflammation of the stomach or intestines, bleeding from surgery, serious injury, or blood donation. Causes of decreased production include iron deficiency, folate deficiency, vitamin B12 deficiency, thalassemia and a number of bone marrow tumors. Causes of increased breakdown include genetic disorders such as sickle cell anemia, infections such as malaria, and certain autoimmune diseases like autoimmune hemolytic anemia. Anemia can also be classified based on the size of the red blood cells and amount of hemoglobin in each cell. If the cells are small, it is called microcytic anemia; if they are large, it is called macrocytic anemia; and if they are normal sized, it is called normocytic anemia. The diagnosis of anemia in men is based on a hemoglobin of less than 130 to 140 g/L (13 to 14 g/dL); in women, it is less than 120 to 130 g/L (12 to 13 g/dL). Further testing is then required to determine the cause. Treatment depends on the specific cause. Certain groups of individuals, such as pregnant women, can benefit from the use of iron pills for prevention. Dietary supplementation, without determining the specific cause, is not recommended. The use of blood transfusions is typically based on a person's signs and symptoms. In those without symptoms, they are not recommended unless hemoglobin levels are less than 60 to 80 g/L (6 to 8 g/dL). These recommendations may also apply to some people with acute bleeding. Erythropoiesis-stimulating agents are only recommended in those with severe anemia. Anemia is the most common blood disorder, affecting about a fifth to a third of the global population. Iron-deficiency anemia is the most common cause of anemia worldwide, and affects nearly one billion people. In 2013, anemia due to iron deficiency resulted in about 183,000 deaths – down from 213,000 deaths in 1990. This condition is most prevalent in children with also an above average prevalence in elderly and women of reproductive age (especially during pregnancy). Anemia is one of the six WHO global nutrition targets for 2025 and for diet-related global targets endorsed by World Health Assembly in 2012 and 2013. Efforts to reach global targets contribute to reaching Sustainable Development Goals (SDGs), with anemia as one of the targets in SDG 2 for achieving zero world hunger. Signs and symptoms A person with anemia may not have any symptoms, depending on the underlying cause, and no symptoms may be noticed, as the anemia is initially mild, and then the symptoms become worse as the anemia worsens. A patient with anemia may report feeling tired, weak, decreased ability to concentrate, and sometimes shortness of breath on exertion. These symptoms are unspecific and none of the symptoms alone or in combination show a good predictive value for the presence of anemia in non-clinical patients. Symptoms of anemia can come on quickly or slowly. Early on there may be few or no symptoms. If the anemia continues slowly (chronic), the body may adapt and compensate for this change. In this case, no symptoms may appear until the anemia becomes more severe. Symptoms can include feeling tired, weak, dizziness, headaches, intolerance to physical exertion, shortness of breath, difficulty concentrating, irregular or rapid heartbeat, cold hands and feet, cold intolerance, pale or yellow skin, poor appetite, easy bruising and bleeding, and muscle weakness. Anemia that develops quickly, often, has more severe symptoms, including, feeling faint, chest pain, sweating, increased thirst, and confusion. There may be also additional symptoms depending on the underlying cause. In more severe anemia, the body may compensate for the lack of oxygen-carrying capability of the blood by increasing cardiac output. The person may have symptoms related to this, such as palpitations, angina (if pre-existing heart disease is present), intermittent claudication of the legs, and symptoms of heart failure. On examination, the signs exhibited may include pallor (pale skin, mucosa, conjunctiva and nail beds), but this is not a reliable sign. Iron-deficiency anemia may give symptoms that can include spoon-shaped nails, restless legs syndrome, and pica (the medical condition indicates the desire for things that are not food, such as ice, dirt, etc.). A blue coloration of the sclera may be noticed in some cases of iron-deficiency anemia. Vitamin B12 deficiency anemia may result in decreased ability to think, memory loss, confusion, personality or mood changes, depression, difficulty walking, blurred vision, and irreversible nerve damage. Other specific causes of anemia may have signs and/or complications such as, jaundice with the rapid break down of red blood cells as with hemolytic anemia, bone abnormalities with thalassemia major, or leg ulcers as seen in sickle cell disease. In severe anemia, there may be signs of a hyperdynamic circulation: tachycardia (a fast heart rate), bounding pulse, flow murmurs, and cardiac ventricular hypertrophy (enlargement). There may be signs of heart failure. Pica, the consumption of non-food items such as ice, paper, wax, grass, hair or dirt, may be a symptom of iron deficiency; although it occurs often in those who have normal levels of hemoglobin. Chronic anemia may result in behavioral disturbances in children as a direct result of impaired neurological development in infants, and reduced academic performance in children of school age. Restless legs syndrome is more common in people with iron-deficiency anemia than in the general population. Causes The causes of anemia may be classified as impaired red blood cell (RBC) production, increased RBC destruction (hemolytic anemia), blood loss and fluid overload (hypervolemia). Several of these may interplay to cause anemia. The most common cause of anemia is blood loss, but this usually does not cause any lasting symptoms unless a relatively impaired RBC production develops, in turn, most commonly by iron deficiency. Impaired production Disturbance of proliferation and differentiation of stem cells Pure red cell aplasia Aplastic anemia affects all kinds of blood cells. Fanconi anemia is a hereditary disorder or defect featuring aplastic anemia and various other abnormalities. Anemia of kidney failure due to insufficient production of the hormone erythropoietin Anemia of endocrine disease Disturbance of proliferation and maturation of erythroblasts Pernicious anemia is a form of megaloblastic anemia due to vitamin B12 deficiency dependent on impaired absorption of vitamin B12. Lack of dietary B12 causes non-pernicious megaloblastic anemia. Anemia of folate deficiency, as with vitamin B12, causes megaloblastic anemia Anemia of prematurity, by diminished erythropoietin response to declining hematocrit levels, combined with blood loss from laboratory testing, generally occurs in premature infants at two to six weeks of age. Iron-deficiency anemia, resulting in deficient heme synthesis Thalassemias, causing deficient globin synthesis Congenital dyserythropoietic anemias, causing ineffective erythropoiesis Anemia of kidney failure (also causing stem cell dysfunction) Other mechanisms of impaired RBC production Myelophthisic anemia or myelophthisis is a severe type of anemia resulting from the replacement of bone marrow by other materials, such as malignant tumors, fibrosis, or granulomas. Myelodysplastic syndrome anemia of chronic inflammation Leukoerythroblastic anemia is caused by space-occupying lesions in the bone marrow that prevent normal production of blood cells. Increased destruction Anemias of increased red blood cell destruction are generally classified as hemolytic anemias. These types generally feature jaundice, and elevated levels of lactate dehydrogenase. Intrinsic (intracorpuscular) abnormalities cause premature destruction. All of these, except paroxysmal nocturnal hemoglobinuria, are hereditary genetic disorders. Hereditary spherocytosis is a hereditary defect that results in defects in the RBC cell membrane, causing the erythrocytes to be sequestered and destroyed by the spleen. Hereditary elliptocytosis is another defect in membrane skeleton proteins. Abetalipoproteinemia, causing defects in membrane lipids Enzyme deficiencies Pyruvate kinase and hexokinase deficiencies, causing defect glycolysis Glucose-6-phosphate dehydrogenase deficiency and glutathione synthetase deficiency, causing increased oxidative stress Hemoglobinopathies Sickle cell anemia Hemoglobinopathies causing unstable hemoglobins Paroxysmal nocturnal hemoglobinuria Extrinsic (extracorpuscular) abnormalities Antibody-mediated Warm autoimmune hemolytic anemia is caused by autoimmune attack against red blood cells, primarily by IgG. It is the most common of the autoimmune hemolytic diseases. It can be idiopathic, that is, without any known cause, drug-associated or secondary to another disease such as systemic lupus erythematosus, or a malignancy, such as chronic lymphocytic leukemia. Cold agglutinin hemolytic anemia is primarily mediated by IgM. It can be idiopathic or result from an underlying condition. Rh disease, one of the causes of hemolytic disease of the newborn Transfusion reaction to blood transfusions Mechanical trauma to red blood cells Microangiopathic hemolytic anemias, including thrombotic thrombocytopenic purpura and disseminated intravascular coagulation Infections, including malaria Heart surgery Haemodialysis Parasitic Trypanosoma congolense alters the surfaces of RBCs of its host and this may explain T. c. induced anemia Blood loss Anemia of prematurity, from frequent blood sampling for laboratory testing, combined with insufficient RBC production Trauma or surgery, causing acute blood loss Gastrointestinal tract lesions, causing either acute bleeds (e.g. variceal lesions, peptic ulcers, hemorrhoids) or chronic blood loss (e.g. angiodysplasia) Gynecologic disturbances, also generally causing chronic blood loss From menstruation, mostly among young women or older women who have fibroids Many type of cancers, including colorectal cancer and cancer of the urinary bladder, may cause acute or chronic blood loss, especially at advanced stages Infection by intestinal nematodes feeding on blood, such as hookworms and the whipworm Trichuris trichiura Iatrogenic anemia, blood loss from repeated blood draws and medical procedures. The roots of the words anemia and ischemia both refer to the basic idea of "lack of blood", but anemia and ischemia are not the same thing in modern medical terminology. The word anemia used alone implies widespread effects from blood that either is too scarce (e.g., blood loss) or is dysfunctional in its oxygen-supplying ability (due to whatever type of hemoglobin or erythrocyte problem). In contrast, the word ischemia refers solely to the lack of blood (poor perfusion). Thus ischemia in a body part can cause localized anemic effects within those tissues. Fluid overload Fluid overload (hypervolemia) causes decreased hemoglobin concentration and apparent anemia: General causes of hypervolemia include excessive sodium or fluid intake, sodium or water retention and fluid shift into the intravascular space. From the 6th week of pregnancy, hormonal changes cause an increase in the mother's blood volume due to an increase in plasma. Intestinal inflammation Certain gastrointestinal disorders can cause anemia. The mechanisms involved are multifactorial and not limited to malabsorption but mainly related to chronic intestinal inflammation, which causes dysregulation of hepcidin that leads to decreased access of iron to the circulation. Helicobacter pylori infection. Gluten-related disorders: untreated celiac disease and non-celiac gluten sensitivity. Anemia can be the only manifestation of celiac disease, in absence of gastrointestinal or any other symptoms. Inflammatory bowel disease. Diagnosis Definitions There are a number of definitions of anemia; reviews provide comparison and contrast of them. A strict but broad definition is an absolute decrease in red blood cell mass, however, a broader definition is a lowered ability of the blood to carry oxygen. An operational definition is a decrease in whole-blood hemoglobin concentration of more than 2 standard deviations below the mean of an age- and sex-matched reference range. It is difficult to directly measure RBC mass, so the hematocrit (amount of RBCs) or the hemoglobin (Hb) in the blood are often used instead to indirectly estimate the value. Hematocrit; however, is concentration dependent and is therefore not completely accurate. For example, during pregnancy a woman's RBC mass is normal but because of an increase in blood volume the hemoglobin and hematocrit are diluted and thus decreased. Another example would be bleeding where the RBC mass would decrease but the concentrations of hemoglobin and hematocrit initially remains normal until fluids shift from other areas of the body to the intravascular space. The anemia is also classified by severity into mild (110 g/L to normal), moderate (80 g/L to 110 g/L), and severe anemia (less than 80 g/L) in adults. Different values are used in pregnancy and children. Testing Anemia is typically diagnosed on a complete blood count. Apart from reporting the number of red blood cells and the hemoglobin level, the automatic counters also measure the size of the red blood cells by flow cytometry, which is an important tool in distinguishing between the causes of anemia. Examination of a stained blood smear using a microscope can also be helpful, and it is sometimes a necessity in regions of the world where automated analysis is less accessible. A blood test will provide counts of white blood cells, red blood cells and platelets. If anemia appears, further tests may determine what type it is, and whether it has a serious cause. although of that, it is possible to refer to the genetic history and physical diagnosis. These tests may also include serum ferritin, iron studies, vitamin B12, genetic testing, and a bone marrow sample, if needed. Reticulocyte counts, and the "kinetic" approach to anemia, have become more common than in the past in the large medical centers of the United States and some other wealthy nations, in part because some automatic counters now have the capacity to include reticulocyte counts. A reticulocyte count is a quantitative measure of the bone marrow's production of new red blood cells. The reticulocyte production index is a calculation of the ratio between the level of anemia and the extent to which the reticulocyte count has risen in response. If the degree of anemia is significant, even a "normal" reticulocyte count actually may reflect an inadequate response. If an automated count is not available, a reticulocyte count can be done manually following special staining of the blood film. In manual examination, activity of the bone marrow can also be gauged qualitatively by subtle changes in the numbers and the morphology of young RBCs by examination under a microscope. Newly formed RBCs are usually slightly larger than older RBCs and show polychromasia. Even where the source of blood loss is obvious, evaluation of erythropoiesis can help assess whether the bone marrow will be able to compensate for the loss and at what rate. When the cause is not obvious, clinicians use other tests, such as: ESR, serum iron, transferrin, RBC folate level, hemoglobin electrophoresis, renal function tests (e.g. serum creatinine) although the tests will depend on the clinical hypothesis that is being investigated. When the diagnosis remains difficult, a bone marrow examination allows direct examination of the precursors to red cells, although is rarely used as is painful, invasive and is hence reserved for cases where severe pathology needs to be determined or excluded. Red blood cell size In the morphological approach, anemia is classified by the size of red blood cells; this is either done automatically or on microscopic examination of a peripheral blood smear. The size is reflected in the mean corpuscular volume (MCV). If the cells are smaller than normal (under 80 fl), the anemia is said to be microcytic; if they are normal size (80–100 fl), normocytic; and if they are larger than normal (over 100 fl), the anemia is classified as macrocytic. This scheme quickly exposes some of the most common causes of anemia; for instance, a microcytic anemia is often the result of iron deficiency. In clinical workup, the MCV will be one of the first pieces of information available, so even among clinicians who consider the "kinetic" approach more useful philosophically, morphology will remain an important element of classification and diagnosis. Limitations of MCV include cases where the underlying cause is due to a combination of factors – such as iron deficiency (a cause of microcytosis) and vitamin B12 deficiency (a cause of macrocytosis) where the net result can be normocytic cells. Production vs. destruction or loss The "kinetic" approach to anemia yields arguably the most clinically relevant classification of anemia. This classification depends on evaluation of several hematological parameters, particularly the blood reticulocyte (precursor of mature RBCs) count. This then yields the classification of defects by decreased RBC production versus increased RBC destruction or loss. Clinical signs of loss or destruction include abnormal peripheral blood smear with signs of hemolysis; elevated LDH suggesting cell destruction; or clinical signs of bleeding, such as guaiac-positive stool, radiographic findings, or frank bleeding. The following is a simplified schematic of this approach: * For instance, sickle cell anemia with superimposed iron deficiency; chronic gastric bleeding with B12 and folate deficiency; and other instances of anemia with more than one cause. ** Confirm by repeating reticulocyte count: ongoing combination of low reticulocyte production index, normal MCV and hemolysis or loss may be seen in bone marrow failure or anemia of chronic disease, with superimposed or related hemolysis or blood loss. Here is a schematic representation of how to consider anemia with MCV as the starting point: Other characteristics visible on the peripheral smear may provide valuable clues about a more specific diagnosis; for example, abnormal white blood cells may point to a cause in the bone marrow. Microcytic Microcytic anemia is primarily a result of hemoglobin synthesis failure/insufficiency, which could be caused by several etiologies: Iron-deficiency anemia is the most common type of anemia overall and it has many causes. RBCs often appear hypochromic (paler than usual) and microcytic (smaller than usual) when viewed with a microscope. Iron-deficiency anemia is due to insufficient dietary intake or absorption of iron to meet the body's needs. Infants, toddlers, and pregnant women have higher than average needs. Increased iron intake is also needed to offset blood losses due to digestive tract issues, frequent blood donations, or heavy menstrual periods. Iron is an essential part of hemoglobin, and low iron levels result in decreased incorporation of hemoglobin into red blood cells. In the United States, 12% of all women of childbearing age have iron deficiency, compared with only 2% of adult men. The incidence is as high as 20% among African American and Mexican American women. In India it is even more than 50%. Studies have linked iron deficiency without anemia to poor school performance and lower IQ in teenage girls, although this may be due to socioeconomic factors. Iron deficiency is the most prevalent deficiency state on a worldwide basis. It is sometimes the cause of abnormal fissuring of the angular (corner) sections of the lips (angular stomatitis). In the United States, the most common cause of iron deficiency is bleeding or blood loss, usually from the gastrointestinal tract. Fecal occult blood testing, upper endoscopy and lower endoscopy should be performed to identify bleeding lesions. In older men and women, the chances are higher that bleeding from the gastrointestinal tract could be due to colon polyps or colorectal cancer. Worldwide, the most common cause of iron-deficiency anemia is parasitic infestation (hookworms, amebiasis, schistosomiasis and whipworms). The Mentzer index (mean cell volume divided by the RBC count) predicts whether microcytic anemia may be due to iron deficiency or thalassemia, although it requires confirmation. Macrocytic Megaloblastic anemia, the most common cause of macrocytic anemia, is due to a deficiency of either vitamin B12, folic acid, or both. Deficiency in folate or vitamin B12 can be due either to inadequate intake or insufficient absorption. Folate deficiency normally does not produce neurological symptoms, while B12 deficiency does. Pernicious anemia is caused by a lack of intrinsic factor, which is required to absorb vitamin B12 from food. A lack of intrinsic factor may arise from an autoimmune condition targeting the parietal cells (atrophic gastritis) that produce intrinsic factor or against intrinsic factor itself. These lead to poor absorption of vitamin B12. Macrocytic anemia can also be caused by the removal of the functional portion of the stomach, such as during gastric bypass surgery, leading to reduced vitamin B12/folate absorption. Therefore, one must always be aware of anemia following this procedure. Hypothyroidism Alcoholism commonly causes a macrocytosis, although not specifically anemia. Other types of liver disease can also cause macrocytosis. Drugs such as methotrexate, zidovudine, and other substances may inhibit DNA replication such as heavy metals Macrocytic anemia can be further divided into "megaloblastic anemia" or "nonmegaloblastic macrocytic anemia". The cause of megaloblastic anemia is primarily a failure of DNA synthesis with preserved RNA synthesis, which results in restricted cell division of the progenitor cells. The megaloblastic anemias often present with neutrophil hypersegmentation (six to 10 lobes). The nonmegaloblastic macrocytic anemias have different etiologies (i.e. unimpaired DNA globin synthesis,) which occur, for example, in alcoholism. In addition to the nonspecific symptoms of anemia, specific features of vitamin B12 deficiency include peripheral neuropathy and subacute combined degeneration of the cord with resulting balance difficulties from posterior column spinal cord pathology. Other features may include a smooth, red tongue and glossitis. The treatment for vitamin B12-deficient anemia was first devised by William Murphy, who bled dogs to make them anemic, and then fed them various substances to see what (if anything) would make them healthy again. He discovered that ingesting large amounts of liver seemed to cure the disease. George Minot and George Whipple then set about to isolate the curative substance chemically and ultimately were able to isolate the vitamin B12 from the liver. All three shared the 1934 Nobel Prize in Medicine. Normocytic Normocytic anemia occurs when the overall hemoglobin levels are decreased, but the red blood cell size (mean corpuscular volume) remains normal. Causes include: Dimorphic A dimorphic appearance on a peripheral blood smear occurs when there are two simultaneous populations of red blood cells, typically of different size and hemoglobin content (this last feature affecting the color of the red blood cell on a stained peripheral blood smear). For example, a person recently transfused for iron deficiency would have small, pale, iron deficient red blood cells (RBCs) and the donor RBCs of normal size and color. Similarly, a person transfused for severe folate or vitamin B12 deficiency would have two cell populations, but, in this case, the patient's RBCs would be larger and paler than the donor's RBCs. A person with sideroblastic anemia (a defect in heme synthesis, commonly caused by alcoholism, but also drugs/toxins, nutritional deficiencies, a few acquired and rare congenital diseases) can have a dimorphic smear from the sideroblastic anemia alone. Evidence for multiple causes appears with an elevated RBC distribution width (RDW), indicating a wider-than-normal range of red cell sizes, also seen in common nutritional anemia. Heinz body anemia Heinz bodies form in the cytoplasm of RBCs and appear as small dark dots under the microscope. In animals, Heinz body anemia has many causes. It may be drug-induced, for example in cats and dogs by acetaminophen (paracetamol), or may be caused by eating various plants or other substances: In cats and dogs after eating either raw or cooked plants from the genus Allium, for example, onions or garlic. In dogs after ingestion of zinc, for example, after eating U.S. pennies minted after 1982. In horses which eat dry or wilted red maple leaves. Hyperanemia Hyperanemia is a severe form of anemia, in which the hematocrit is below 10%. Refractory anemia Refractory anemia, an anemia which does not respond to treatment, is often seen secondary to myelodysplastic syndromes. Iron-deficiency anemia may also be refractory as a manifestation of gastrointestinal problems which disrupt iron absorption or cause occult bleeding. Transfusion dependent Transfusion dependent anemia is a form of anemia where ongoing blood transfusion are required. Most people with myelodysplastic syndrome develop this state at some point in time. Beta thalassemia may also result in transfusion dependence. Concerns from repeated blood transfusions include iron overload. This iron overload may require chelation therapy. Treatment The global market for anemia treatments is estimated at more than USD 23 billion per year and is fast growing because of the rising prevalence and awareness of anemia. The types of anemia treated with drugs are iron-deficiency anemia, thalassemia, aplastic anemia, hemolytic anemia, sickle cell anemia, and pernicious anemia, the most important of them being deficiency and sickle cell anemia with together 60% of market share because of highest prevalence as well as higher treatment costs compared with other types. Treatment for anemia depends on cause and severity. Vitamin supplements given orally (folic acid or vitamin B12) or intramuscularly (vitamin B12) will replace specific deficiencies. Apart from that, iron supplements, antibiotics, immunosuppressant, bone marrow stimulants, corticosteroids, gene therapy and iron chelating agents are forms of anemia treatment drugs, with immunosuppressants and corticosteroids accounting for 58% of the market share. A paradigm shift towards gene therapy and monoclonal antibody therapies is observed. Oral iron Nutritional iron deficiency is common in developing nations. An estimated two-thirds of children and of women of childbearing age in most developing nations are estimated to have iron deficiency without anemia with one-third of them having an iron deficiency with anemia. Iron deficiency due to inadequate dietary iron intake is rare in men and postmenopausal women. The diagnosis of iron deficiency mandates a search for potential sources of blood loss, such as gastrointestinal bleeding from ulcers or colon cancer. Mild to moderate iron-deficiency anemia is treated by oral iron supplementation with ferrous sulfate, ferrous fumarate, or ferrous gluconate. Daily iron supplements have been shown to be effective in reducing anemia in women of childbearing age. When taking iron supplements, stomach upset or darkening of the feces are commonly experienced. The stomach upset can be alleviated by taking the iron with food; however, this decreases the amount of iron absorbed. Vitamin C aids in the body's ability to absorb iron, so taking oral iron supplements with orange juice is of benefit. In the anemia of chronic kidney disease, recombinant erythropoietin or epoetin alfa is recommended to stimulate RBC production, and if iron deficiency and inflammation are also present, concurrent parenteral iron is also recommended. Injectable iron In cases where oral iron has either proven ineffective, would be too slow (for example, pre-operatively), or where absorption is impeded (for example in cases of inflammation), parenteral iron preparations can be used. Parenteral iron can improve iron stores rapidly and is also effective for treating people with postpartum haemorrhage, inflammatory bowel disease, and chronic heart failure. The body can absorb up to 6 mg iron daily from the gastrointestinal tract. In many cases, the patient has a deficit of over 1,000 mg of iron which would require several months to replace. This can be given concurrently with erythropoietin to ensure sufficient iron for increased rates of erythropoiesis. Blood transfusions Blood transfusions in those without symptoms is not recommended until the hemoglobin is below 60 to 80 g/L (6 to 8 g/dL). In those with coronary artery disease who are not actively bleeding transfusions are only recommended when the hemoglobin is below 70 to 80g/L (7 to 8 g/dL). Transfusing earlier does not improve survival. Transfusions otherwise should only be undertaken in cases of cardiovascular instability. A 2012 review concluded that when considering blood transfusions for anaemia in people with advanced cancer who have fatigue and breathlessness (not related to cancer treatment or haemorrhage), consideration should be given to whether there are alternative strategies can be tried before a blood transfusion. Vitamin B12 intramuscular injections In many cases, vitamin B12 is used by intramuscular injection in severe cases or cases of malabsorption of dietary-B12. Pernicious anemia caused by loss of intrinsic factor cannot be prevented. If there are other, reversible causes of low vitamin B12 levels, the cause must be treated. Vitamin B12 deficiency anemia is usually easily treated by providing the necessary level of vitamin B12 supplementation. The injections are quick-acting, and symptoms usually go away within one to two weeks. As the condition improves, doses are reduced to weeks and then can be given monthly. Intramuscular therapy leads to more rapid improvement and should be considered in patients with severe deficiency or severe neurologic symptoms. Treatment should begin rapidly for severe neurological symptoms, as some changes can become permanent. In some individuals lifelong treatment may be needed. Erythropoiesis-stimulating agents The objective for the administration of an erythropoiesis-stimulating agent (ESA) is to maintain hemoglobin at the lowest level that both minimizes transfusions and meets the individual person's needs. They should not be used for mild or moderate anemia. They are not recommended in people with chronic kidney disease unless hemoglobin levels are less than 10 g/dL or they have symptoms of anemia. Their use should be along with parenteral iron. The 2020 Cochrane Anaesthesia Review Group review of erythropoietin (EPO) plus iron versus control treatment including placebo or iron for preoperative anaemic adults undergoing non-cardiac surgery demonstrated that patients were much less likely to require red cell transfusion and in those transfused, the volumes were unchanged (mean difference -0.09, 95% CI -0.23 to 0.05). Pre-operative hemoglobin concentration was increased in those receiving 'high dose' EPO, but not 'low dose'. Hyperbaric oxygen Treatment of exceptional blood loss (anemia) is recognized as an indication for hyperbaric oxygen (HBO) by the Undersea and Hyperbaric Medical Society. The use of HBO is indicated when oxygen delivery to tissue is not sufficient in patients who cannot be given blood transfusions for medical or religious reasons. HBO may be used for medical reasons when threat of blood product incompatibility or concern for transmissible disease are factors. The beliefs of some religions (ex: Jehovah's Witnesses) may require they use the HBO method. A 2005 review of the use of HBO in severe anemia found all publications reported positive results. Preoperative anemia An estimated 30% of adults who require non-cardiac surgery have anemia. In order to determine an appropriate preoperative treatment, it is suggested that the cause of anemia be first determined. There is moderate level medical evidence that supports a combination of iron supplementation and erythropoietin treatment to help reduce the requirement for red blood cell transfusions after surgery in those who have preoperative anemia. Epidemiology Anemia affects 27% of the world's population with iron-deficiency anemia accounting for more than 60% of it. A moderate degree of iron-deficiency anemia affected approximately 610 million people worldwide or 8.8% of the population. It is somewhat more common in females (9.9%) than males (7.8%). Mild iron-deficiency anemia affects another 375 million. Severe anaemia is prevalent globally, and especially in sub-Saharan Africa where it is associated with infections including malaria and invasive bacterial infections. History Signs of severe anemia in human bones from 4000 years ago have been uncovered in Thailand.	Biology and health sciences	Non-infectious disease	null
84228	https://en.wikipedia.org/wiki/Compound%20bow	Compound bow	In modern archery, a compound bow is a bow that uses a levering system, usually of cables and pulleys, to bend the limbs. The compound bow was first developed in 1966 by Holless Wilbur Allen in North Kansas City, Missouri, and a US patent was granted in 1969. Compound bows are widely used in target practice and hunting. Compound bows are typically constructed of man-made materials such as fiberglass and carbon fiber, while traditional bows and warbows usually are entirely or partially made of wood or bamboo. The pulley/cam system grants the user a mechanical advantage, and so the limbs of a compound bow are much stiffer than those of a recurve bow or longbow. This rigidity makes the compound bow more energy-efficient than traditional bows, as less energy is dissipated in limb movement. The higher-rigidity, more advanced construction also improves accuracy by reducing the bow's sensitivity to changes in temperature and humidity. In literature of the early 20th century, before the invention of compound bows, composite bows were described as "compound". Construction A bow's central mount for other components (limbs, sights, stabilizers and quivers) is called the riser. Risers are designed to be as rigid as possible. The central riser of a compound bow is usually made of aluminum, magnesium alloy, or carbon fiber and many are made of 7075 aluminum alloy. Limbs are made of fiberglass-based composite materials, or occasionally wood, and able to withstand high tensile and compressive forces. The limbs store the kinetic energy of the bow – no energy is stored in the pulleys and cables. Draw weights of adult compound bows range is between , which can create arrow speeds of . In the most common configuration, there is a cam or wheel at the end of each limb. The shape of the cam may vary somewhat between different bow designs. There are several different concepts of using the cams to store energy in the limbs, and these all fall under a category called bow eccentrics. The four most common types of bow eccentrics are Single Cam, Hybrid Cam, Dual Cam and Binary Cam. However, there are also other less common designs, like the Quad Cam and Hinged. Cams are often described using their "let-off" rating. As a cam is rotated, the force required to hold the bow in position reaches a peak and then decreases as the bow approaches maximum extension (a position known as "the wall"). The percent-difference between the maximum force encountered during the draw and the force required to hold the bow in full extension is the "let-off". This value is commonly between 65% and 80% of the peak weight for recently designed compound bows, although some older compound bows provided a let-off of only 50% and some recent designs achieve let-offs in excess of 90%. As the string is drawn the cam turns and imparts force to compress the limb. Initially, the archer has the 'short' side of the cam, with the leverage being a mechanical disadvantage. High energy input is therefore required. When near full draw is reached, the cam has turned to its full extent, the archer has gained mechanical advantage, and the least amount of force needs to be applied to the string to keep the limbs bent. This is known as "let off". The lower holding weight enables the archer to maintain the bow fully drawn and take more time to aim. This let-off enables the archer to accurately shoot a compound bow with a much higher peak draw weight than other bows (see below). However, there are some youth-oriented compound bows with low draw weights that have no let-off and have a maximum draw length deliberately set farther than the majority of young shooters would reach. This effectively makes the bow function very similar to a recurve, with the draw length determined by the shooter's preferred anchor point. This removes the necessity to adjust the bow draw length or use a different bow for different shooters (or to change bows as the shooter gets older). An example of this type of bow is the Genesis, which is standard equipment in the U.S. National Archery in the Schools Program. Compound bow strings and cables are normally made of high-modulus polyethylene and are designed to have great tensile strength and minimal stretchability, so that the bow transfers its energy to the arrow as efficiently and durably as possible. In earlier models of compound bows, the cables were often made of plastic-coated steel. Comparison to other bow types Technical advantages The function of the cam systems (known as the 'eccentrics') is to maximize the energy storage throughout the draw cycle and provide let-off at the end of the cycle (less holding weight at full draw). A traditional recurve bow has a very linear draw weight curve - meaning that as the bow is drawn back, the draw force becomes heavier with each inch of draw (and most difficult at full draw). Therefore, little energy is stored in the first half of the draw, and much more energy at the end where the draw weight is heaviest. The compound bow operates with a different weight profile, reaching its peak weight within the first few inches of the draw, and remaining more flat and constant until the end of the cycle where the cams "let-off" and allow a reduced holding weight. This manipulation of the peak weight throughout the draw (accomplished by the elliptical shape of the cams that change leverage and mechanical advantage) is why compound bows store more energy and shoot faster than an equivalent peak weight recurve bow or longbow. The design of the cams directly controls the acceleration of the arrow. What is termed a "soft cam" will accelerate the arrow more gently than a "harder" cam. Novice archers will typically shoot a soft cam whereas a more advanced archer may choose to use a harder cam to gain speed. Bows can be had with a variety of cams, in a full spectrum from soft to hard. Some pulley systems use a single cam at the bottom of the bow and a round idler wheel at the top of the bow instead of two identical cams. This design eliminates the need for a separate control cable and instead uses a single long string that begins at the cam on the bottom of the bow, travels over the wheel on top, and back to the bottom cam. A separate buss cable then connects the bottom cam to the top limb. When a compound bow is drawn, the limbs are pulled in toward each other by the cables, unlike a longbow or recurve where the limbs flex in the direction of the bow string. This difference allows modern compounds to have limbs that are closer to horizontal instead of angled. The horizontal, or "parallel" limb configuration minimizes the recoil and vibration felt by the shooter when the arrow is released, as the forces going upward at the top limb and downward at the bottom limb cancel each other out. The pulley system will usually include some rubber-covered blocks that act as draw-stops. These provide a solid "wall" that the archer can draw against. These draw stops can be adjusted to suit the archer's optimum draw-length, which helps the archer achieve a consistent anchor point and a consistent amount of force imparted to the arrow on every shot, further increasing accuracy. Technical disadvantages The relatively larger number of moving parts requires additional maintenance and creates more potential points of failure. Dry firing is more likely to damage or destroy a compound bow due to the greater amount of energy stored and released. Unlike traditional bows, draw length and let-off adjustments as well as string or cable replacements often require a bow press, a specialized tool used for compressing the limbs to take tension off the cables and string. Drawing a compound bow with just the fingers increases the likelihood of torquing the bowstring, which could derail it from the cams. The use of a mechanical release-aid is often required to avoid this. Usually heavier than recurves and longbows. Circumstantial advantages Compound archers often use a mechanical release aid to hold and release the string. This attaches to the bowstring near the point where the arrow attaches, the nocking point, and permits the archer to release the string with a squeeze of a trigger or a slight increase of tension. The use of a release aid gives a more consistent release than the use of fingers on the string as it minimises the arrow oscillation which is inevitable when the bowstring is released directly from the fingers. In tournaments, competition rules for compound archers allow bows with a sighting system, consisting of a "peep sight" held within the bowstring that acts as a back sight, however front sights attached to the riser are allowable in other classes. Some front sights are magnifying and/or adjustable for targets at different distances. Some sights have multiple "pins" set up for targets at different distances. Circumstantial disadvantages The relatively low holding weight of a compound bow compared to a recurve bow makes the compound more sensitive to certain shooting form faults when the archer is at full draw. In particular, it's easier for the archer to torque (twist) the bow around the vertical axis, leading to left-right errors, and also a plucked or snatched release can have more effect. Specifications AMO (Archery Manufacturers and merchants Organization, the former name of the body now known as the Archery Trade Association) standard draw length is the distance from the string at full draw to the lowest point on the grip plus . Because the draw force may increase more or less rapidly, and again drop off more or less rapidly when approaching peak draw, bows of the same peak draw force can store different amounts of energy. Norbert Mullaney has defined the ratio of stored energy to peak draw force (S.E./P.D.F.). This is usually around but can reach . The efficiency of bows also varies. Normally between 70 and 85% of the stored energy is transferred to the arrow. This stored energy is referred to as potential energy. When transferred to the arrow it is referred to as kinetic energy. The product of S.E./P.D.F. and efficiency can be called the power factor. There are two measurement standards of this quantity – ATA and IBO speed. ATA is defined as the initial velocity of a arrow when shot from a bow with a peak draw weight of and draw length of . IBO speed is defined as the initial velocity of an arrow with a weight of per pound of draw weight. While many manufacturers measure IBO speeds using a draw weight of and draw length of , the IBO standard allows a draw weight of as high as , and does not specify a draw length. The average IBO speed for the majority of compound bows on the market hovers around 310–320 feet per second. Brace height is the distance from the pivot point of the grip to the string at rest. Typically a shorter brace height will result in an increased power stroke, but comes at the price of a bow that's less forgiving to shooter error and having harsher string slap. Arrows Arrows used with compound bows do not differ significantly from those used with recurve bows, being typically either aluminum alloy, carbon fiber, or a composite of the two materials. Wooden arrows are not commonly used on compound bows because of their fragility. Most arrows in use today are of the carbon fiber variety. An important distinction arrow-wise between recurve bows and compound bows is that of arrow spine. Compound bows and target recurve bows with fully center-shot cutaway risers tend to be very forgiving in regard to spine selection. Modern compound bows are typically equipped with substantially stiffer arrows than an equivalent draw-length and draw-weight recurve bow would be. Another advantage of the center-shot riser is that the arrow need not bend around the riser (nearly as much or at all) during the shot. Fine-tuning may be accomplished by adjustment of the arrow rest, or nock point on the string, rather than by changing arrow-length and tip weight. Manufacturers produce arrow shafts with different weights, different spines (stiffness), and different lengths in the same model of shaft to accommodate different draw weights and lengths, matched to archers' different styles, preferences and physical attributes. Arrow stiffness (spine) is an important parameter in finding arrows that will shoot accurately from any particular bow (see Archer's paradox), the spine varying with both the construction and length of the arrow. Another important consideration is that the IBO (International Bowhunting Organization) recommends at least of draw weight as a safety buffer. This means a bow that draws would need at least a finished-with-tip arrow. Shooting arrows lighter than this guideline risks damage to the bow similar to that caused by dry-firing, which can in turn cause injury to the archer or anyone standing nearby. Shooting arrows that are too light also voids most manufacturer warranties.	Technology	Archery	null
84400	https://en.wikipedia.org/wiki/Zero-point%20energy	Zero-point energy	Zero-point energy (ZPE) is the lowest possible energy that a quantum mechanical system may have. Unlike in classical mechanics, quantum systems constantly fluctuate in their lowest energy state as described by the Heisenberg uncertainty principle. Therefore, even at absolute zero, atoms and molecules retain some vibrational motion. Apart from atoms and molecules, the empty space of the vacuum also has these properties. According to quantum field theory, the universe can be thought of not as isolated particles but continuous fluctuating fields: matter fields, whose quanta are fermions (i.e., leptons and quarks), and force fields, whose quanta are bosons (e.g., photons and gluons). All these fields have zero-point energy. These fluctuating zero-point fields lead to a kind of reintroduction of an aether in physics since some systems can detect the existence of this energy. However, this aether cannot be thought of as a physical medium if it is to be Lorentz invariant such that there is no contradiction with Einstein's theory of special relativity. The notion of a zero-point energy is also important for cosmology, and physics currently lacks a full theoretical model for understanding zero-point energy in this context; in particular, the discrepancy between theorized and observed vacuum energy in the universe is a source of major contention. Yet according to Einstein's theory of general relativity, any such energy would gravitate, and the experimental evidence from the expansion of the universe, dark energy and the Casimir effect shows any such energy to be exceptionally weak. One proposal that attempts to address this issue is to say that the fermion field has a negative zero-point energy, while the boson field has positive zero-point energy and thus these energies somehow cancel out each other. This idea would be true if supersymmetry were an exact symmetry of nature; however, the Large Hadron Collider at CERN has so far found no evidence to support it. Moreover, it is known that if supersymmetry is valid at all, it is at most a broken symmetry, only true at very high energies, and no one has been able to show a theory where zero-point cancellations occur in the low-energy universe we observe today. This discrepancy is known as the cosmological constant problem and it is one of the greatest unsolved mysteries in physics. Many physicists believe that "the vacuum holds the key to a full understanding of nature". Etymology and terminology The term zero-point energy (ZPE) is a translation from the German . Sometimes used interchangeably with it are the terms zero-point radiation and ground state energy. The term zero-point field (ZPF) can be used when referring to a specific vacuum field, for instance the QED vacuum which specifically deals with quantum electrodynamics (e.g., electromagnetic interactions between photons, electrons and the vacuum) or the QCD vacuum which deals with quantum chromodynamics (e.g., color charge interactions between quarks, gluons and the vacuum). A vacuum can be viewed not as empty space but as the combination of all zero-point fields. In quantum field theory this combination of fields is called the vacuum state, its associated zero-point energy is called the vacuum energy and the average energy value is called the vacuum expectation value (VEV) also called its condensate. Overview In classical mechanics all particles can be thought of as having some energy made up of their potential energy and kinetic energy. Temperature, for example, arises from the intensity of random particle motion caused by kinetic energy (known as Brownian motion). As temperature is reduced to absolute zero, it might be thought that all motion ceases and particles come completely to rest. In fact, however, kinetic energy is retained by particles even at the lowest possible temperature. The random motion corresponding to this zero-point energy never vanishes; it is a consequence of the uncertainty principle of quantum mechanics. The uncertainty principle states that no object can ever have precise values of position and velocity simultaneously. The total energy of a quantum mechanical object (potential and kinetic) is described by its Hamiltonian which also describes the system as a harmonic oscillator, or wave function, that fluctuates between various energy states (see wave-particle duality). All quantum mechanical systems undergo fluctuations even in their ground state, a consequence of their wave-like nature. The uncertainty principle requires every quantum mechanical system to have a fluctuating zero-point energy greater than the minimum of its classical potential well. This results in motion even at absolute zero. For example, liquid helium does not freeze under atmospheric pressure regardless of temperature due to its zero-point energy. Given the equivalence of mass and energy expressed by Albert Einstein's , any point in space that contains energy can be thought of as having mass to create particles. Modern physics has developed quantum field theory (QFT) to understand the fundamental interactions between matter and forces; it treats every single point of space as a quantum harmonic oscillator. According to QFT the universe is made up of matter fields, whose quanta are fermions (i.e. leptons and quarks), and force fields, whose quanta are bosons (e.g. photons and gluons). All these fields have zero-point energy. Recent experiments support the idea that particles themselves can be thought of as excited states of the underlying quantum vacuum, and that all properties of matter are merely vacuum fluctuations arising from interactions of the zero-point field. The idea that "empty" space can have an intrinsic energy associated with it, and that there is no such thing as a "true vacuum" is seemingly unintuitive. It is often argued that the entire universe is completely bathed in the zero-point radiation, and as such it can add only some constant amount to calculations. Physical measurements will therefore reveal only deviations from this value. For many practical calculations zero-point energy is dismissed by fiat in the mathematical model as a term that has no physical effect. Such treatment causes problems however, as in Einstein's theory of general relativity the absolute energy value of space is not an arbitrary constant and gives rise to the cosmological constant. For decades most physicists assumed that there was some undiscovered fundamental principle that will remove the infinite zero-point energy and make it completely vanish. If the vacuum has no intrinsic, absolute value of energy it will not gravitate. It was believed that as the universe expands from the aftermath of the Big Bang, the energy contained in any unit of empty space will decrease as the total energy spreads out to fill the volume of the universe; galaxies and all matter in the universe should begin to decelerate. This possibility was ruled out in 1998 by the discovery that the expansion of the universe is not slowing down but is in fact accelerating, meaning empty space does indeed have some intrinsic energy. The discovery of dark energy is best explained by zero-point energy, though it still remains a mystery as to why the value appears to be so small compared to the huge value obtained through theory – the cosmological constant problem. Many physical effects attributed to zero-point energy have been experimentally verified, such as spontaneous emission, Casimir force, Lamb shift, magnetic moment of the electron and Delbrück scattering. These effects are usually called "radiative corrections". In more complex nonlinear theories (e.g. QCD) zero-point energy can give rise to a variety of complex phenomena such as multiple stable states, symmetry breaking, chaos and emergence. Active areas of research include the effects of virtual particles, quantum entanglement, the difference (if any) between inertial and gravitational mass, variation in the speed of light, a reason for the observed value of the cosmological constant and the nature of dark energy. History Early aether theories Zero-point energy evolved from historical ideas about the vacuum. To Aristotle the vacuum was , "the empty"; i.e., space independent of body. He believed this concept violated basic physical principles and asserted that the elements of fire, air, earth, and water were not made of atoms, but were continuous. To the atomists the concept of emptiness had absolute character: it was the distinction between existence and nonexistence. Debate about the characteristics of the vacuum were largely confined to the realm of philosophy, it was not until much later on with the beginning of the renaissance, that Otto von Guericke invented the first vacuum pump and the first testable scientific ideas began to emerge. It was thought that a totally empty volume of space could be created by simply removing all gases. This was the first generally accepted concept of the vacuum. Late in the 19th century, however, it became apparent that the evacuated region still contained thermal radiation. The existence of the aether as a substitute for a true void was the most prevalent theory of the time. According to the successful electromagnetic aether theory based upon Maxwell's electrodynamics, this all-encompassing aether was endowed with energy and hence very different from nothingness. The fact that electromagnetic and gravitational phenomena were transmitted in empty space was considered evidence that their associated aethers were part of the fabric of space itself. However Maxwell noted that for the most part these aethers were ad hoc: Moreever, the results of the Michelson–Morley experiment in 1887 were the first strong evidence that the then-prevalent aether theories were seriously flawed, and initiated a line of research that eventually led to special relativity, which ruled out the idea of a stationary aether altogether. To scientists of the period, it seemed that a true vacuum in space might be created by cooling and thus eliminating all radiation or energy. From this idea evolved the second concept of achieving a real vacuum: cool a region of space down to absolute zero temperature after evacuation. Absolute zero was technically impossible to achieve in the 19th century, so the debate remained unsolved. Second quantum theory In 1900, Max Planck derived the average energy of a single energy radiator, e.g., a vibrating atomic unit, as a function of absolute temperature: where is the Planck constant, is the frequency, is the Boltzmann constant, and is the absolute temperature. The zero-point energy makes no contribution to Planck's original law, as its existence was unknown to Planck in 1900. The concept of zero-point energy was developed by Max Planck in Germany in 1911 as a corrective term added to a zero-grounded formula developed in his original quantum theory in 1900. In 1912, Max Planck published the first journal article to describe the discontinuous emission of radiation, based on the discrete quanta of energy. In Planck's "second quantum theory" resonators absorbed energy continuously, but emitted energy in discrete energy quanta only when they reached the boundaries of finite cells in phase space, where their energies became integer multiples of . This theory led Planck to his new radiation law, but in this version energy resonators possessed a zero-point energy, the smallest average energy a resonator could take on. Planck's radiation equation contained a residual energy factor, one , as an additional term dependent on the frequency , which was greater than zero (where is the Planck constant). It is therefore widely agreed that "Planck's equation marked the birth of the concept of zero-point energy." In a series of papers from 1911 to 1913, Planck found the average energy of an oscillator to be: Soon, the idea of zero-point energy attracted the attention of Albert Einstein and his assistant Otto Stern. In 1913 they published a paper that attempted to prove the existence of zero-point energy by calculating the specific heat of hydrogen gas and compared it with the experimental data. However, after assuming they had succeeded, they retracted support for the idea shortly after publication because they found Planck's second theory may not apply to their example. In a letter to Paul Ehrenfest of the same year Einstein declared zero-point energy "dead as a doornail". Zero-point energy was also invoked by Peter Debye, who noted that zero-point energy of the atoms of a crystal lattice would cause a reduction in the intensity of the diffracted radiation in X-ray diffraction even as the temperature approached absolute zero. In 1916 Walther Nernst proposed that empty space was filled with zero-point electromagnetic radiation. With the development of general relativity Einstein found the energy density of the vacuum to contribute towards a cosmological constant in order to obtain static solutions to his field equations; the idea that empty space, or the vacuum, could have some intrinsic energy associated with it had returned, with Einstein stating in 1920: and Francis Simon (1923), who worked at Walther Nernst's laboratory in Berlin, studied the melting process of chemicals at low temperatures. Their calculations of the melting points of hydrogen, argon and mercury led them to conclude that the results provided evidence for a zero-point energy. Moreover, they suggested correctly, as was later verified by Simon (1934), that this quantity was responsible for the difficulty in solidifying helium even at absolute zero. In 1924 Robert Mulliken provided direct evidence for the zero-point energy of molecular vibrations by comparing the band spectrum of 10BO and 11BO: the isotopic difference in the transition frequencies between the ground vibrational states of two different electronic levels would vanish if there were no zero-point energy, in contrast to the observed spectra. Then just a year later in 1925, with the development of matrix mechanics in Werner Heisenberg's article "Quantum theoretical re-interpretation of kinematic and mechanical relations" the zero-point energy was derived from quantum mechanics. In 1913 Niels Bohr had proposed what is now called the Bohr model of the atom, but despite this it remained a mystery as to why electrons do not fall into their nuclei. According to classical ideas, the fact that an accelerating charge loses energy by radiating implied that an electron should spiral into the nucleus and that atoms should not be stable. This problem of classical mechanics was nicely summarized by James Hopwood Jeans in 1915: "There would be a very real difficulty in supposing that the (force) law held down to the zero values of . For the forces between two charges at zero distance would be infinite; we should have charges of opposite sign continually rushing together and, when once together, no force would tend to shrink into nothing or to diminish indefinitely in size." The resolution to this puzzle came in 1926 when Erwin Schrödinger introduced the Schrödinger equation. This equation explained the new, non-classical fact that an electron confined to be close to a nucleus would necessarily have a large kinetic energy so that the minimum total energy (kinetic plus potential) actually occurs at some positive separation rather than at zero separation; in other words, zero-point energy is essential for atomic stability. Quantum field theory and beyond In 1926, Pascual Jordan published the first attempt to quantize the electromagnetic field. In a joint paper with Max Born and Werner Heisenberg he considered the field inside a cavity as a superposition of quantum harmonic oscillators. In his calculation he found that in addition to the "thermal energy" of the oscillators there also had to exist an infinite zero-point energy term. He was able to obtain the same fluctuation formula that Einstein had obtained in 1909. However, Jordan did not think that his infinite zero-point energy term was "real", writing to Einstein that "it is just a quantity of the calculation having no direct physical meaning". Jordan found a way to get rid of the infinite term, publishing a joint work with Pauli in 1928, performing what has been called "the first infinite subtraction, or renormalisation, in quantum field theory". Building on the work of Heisenberg and others, Paul Dirac's theory of emission and absorption (1927) was the first application of the quantum theory of radiation. Dirac's work was seen as crucially important to the emerging field of quantum mechanics; it dealt directly with the process in which "particles" are actually created: spontaneous emission. Dirac described the quantization of the electromagnetic field as an ensemble of harmonic oscillators with the introduction of the concept of creation and annihilation operators of particles. The theory showed that spontaneous emission depends upon the zero-point energy fluctuations of the electromagnetic field in order to get started. In a process in which a photon is annihilated (absorbed), the photon can be thought of as making a transition into the vacuum state. Similarly, when a photon is created (emitted), it is occasionally useful to imagine that the photon has made a transition out of the vacuum state. In the words of Dirac: Contemporary physicists, when asked to give a physical explanation for spontaneous emission, generally invoke the zero-point energy of the electromagnetic field. This view was popularized by Victor Weisskopf who in 1935 wrote: This view was also later supported by Theodore Welton (1948), who argued that spontaneous emission "can be thought of as forced emission taking place under the action of the fluctuating field". This new theory, which Dirac coined quantum electrodynamics (QED), predicted a fluctuating zero-point or "vacuum" field existing even in the absence of sources. Throughout the 1940s improvements in microwave technology made it possible to take more precise measurements of the shift of the levels of a hydrogen atom, now known as the Lamb shift, and measurement of the magnetic moment of the electron. Discrepancies between these experiments and Dirac's theory led to the idea of incorporating renormalisation into QED to deal with zero-point infinities. Renormalization was originally developed by Hans Kramers and also Victor Weisskopf (1936), and first successfully applied to calculate a finite value for the Lamb shift by Hans Bethe (1947). As per spontaneous emission, these effects can in part be understood with interactions with the zero-point field. But in light of renormalisation being able to remove some zero-point infinities from calculations, not all physicists were comfortable attributing zero-point energy any physical meaning, viewing it instead as a mathematical artifact that might one day be eliminated. In Wolfgang Pauli's 1945 Nobel lecture he made clear his opposition to the idea of zero-point energy stating "It is clear that this zero-point energy has no physical reality". In 1948 Hendrik Casimir showed that one consequence of the zero-point field is an attractive force between two uncharged, perfectly conducting parallel plates, the so-called Casimir effect. At the time, Casimir was studying the properties of colloidal solutions. These are viscous materials, such as paint and mayonnaise, that contain micron-sized particles in a liquid matrix. The properties of such solutions are determined by Van der Waals forces – short-range, attractive forces that exist between neutral atoms and molecules. One of Casimir's colleagues, Theo Overbeek, realized that the theory that was used at the time to explain Van der Waals forces, which had been developed by Fritz London in 1930, did not properly explain the experimental measurements on colloids. Overbeek therefore asked Casimir to investigate the problem. Working with Dirk Polder, Casimir discovered that the interaction between two neutral molecules could be correctly described only if the fact that light travels at a finite speed was taken into account. Soon afterwards after a conversation with Bohr about zero-point energy, Casimir noticed that this result could be interpreted in terms of vacuum fluctuations. He then asked himself what would happen if there were two mirrors – rather than two molecules – facing each other in a vacuum. It was this work that led to his prediction of an attractive force between reflecting plates. The work by Casimir and Polder opened up the way to a unified theory of van der Waals and Casimir forces and a smooth continuum between the two phenomena. This was done by Lifshitz (1956) in the case of plane parallel dielectric plates. The generic name for both van der Waals and Casimir forces is dispersion forces, because both of them are caused by dispersions of the operator of the dipole moment. The role of relativistic forces becomes dominant at orders of a hundred nanometers. In 1951 Herbert Callen and Theodore Welton proved the quantum fluctuation-dissipation theorem (FDT) which was originally formulated in classical form by Nyquist (1928) as an explanation for observed Johnson noise in electric circuits. The fluctuation-dissipation theorem showed that when something dissipates energy, in an effectively irreversible way, a connected heat bath must also fluctuate. The fluctuations and the dissipation go hand in hand; it is impossible to have one without the other. The implication of FDT being that the vacuum could be treated as a heat bath coupled to a dissipative force and as such energy could, in part, be extracted from the vacuum for potentially useful work. FDT has been shown to be true experimentally under certain quantum, non-classical, conditions. In 1963 the Jaynes–Cummings model was developed describing the system of a two-level atom interacting with a quantized field mode (i.e. the vacuum) within an optical cavity. It gave nonintuitive predictions such as that an atom's spontaneous emission could be driven by field of effectively constant frequency (Rabi frequency). In the 1970s experiments were being performed to test aspects of quantum optics and showed that the rate of spontaneous emission of an atom could be controlled using reflecting surfaces. These results were at first regarded with suspicion in some quarters: it was argued that no modification of a spontaneous emission rate would be possible, after all, how can the emission of a photon be affected by an atom's environment when the atom can only "see" its environment by emitting a photon in the first place? These experiments gave rise to cavity quantum electrodynamics (CQED), the study of effects of mirrors and cavities on radiative corrections. Spontaneous emission can be suppressed (or "inhibited") or amplified. Amplification was first predicted by Purcell in 1946 (the Purcell effect) and has been experimentally verified. This phenomenon can be understood, partly, in terms of the action of the vacuum field on the atom. Uncertainty principle Zero-point energy is fundamentally related to the Heisenberg uncertainty principle. Roughly speaking, the uncertainty principle states that complementary variables (such as a particle's position and momentum, or a field's value and derivative at a point in space) cannot simultaneously be specified precisely by any given quantum state. In particular, there cannot exist a state in which the system simply sits motionless at the bottom of its potential well, for then its position and momentum would both be completely determined to arbitrarily great precision. Therefore, the lowest-energy state (the ground state) of the system must have a distribution in position and momentum that satisfies the uncertainty principle, which implies its energy must be greater than the minimum of the potential well. Near the bottom of a potential well, the Hamiltonian of a general system (the quantum-mechanical operator giving its energy) can be approximated as a quantum harmonic oscillator, where is the minimum of the classical potential well. The uncertainty principle tells us that making the expectation values of the kinetic and potential terms above satisfy The expectation value of the energy must therefore be at least where is the angular frequency at which the system oscillates. A more thorough treatment, showing that the energy of the ground state actually saturates this bound and is exactly , requires solving for the ground state of the system. Atomic physics The idea of a quantum harmonic oscillator and its associated energy can apply to either an atom or a subatomic particle. In ordinary atomic physics, the zero-point energy is the energy associated with the ground state of the system. The professional physics literature tends to measure frequency, as denoted by above, using angular frequency, denoted with and defined by . This leads to a convention of writing the Planck constant with a bar through its top () to denote the quantity . In these terms, an example of zero-point energy is the above associated with the ground state of the quantum harmonic oscillator. In quantum mechanical terms, the zero-point energy is the expectation value of the Hamiltonian of the system in the ground state. If more than one ground state exists, they are said to be degenerate. Many systems have degenerate ground states. Degeneracy occurs whenever there exists a unitary operator which acts non-trivially on a ground state and commutes with the Hamiltonian of the system. According to the third law of thermodynamics, a system at absolute zero temperature exists in its ground state; thus, its entropy is determined by the degeneracy of the ground state. Many systems, such as a perfect crystal lattice, have a unique ground state and therefore have zero entropy at absolute zero. It is also possible for the highest excited state to have absolute zero temperature for systems that exhibit negative temperature. The wave function of the ground state of a particle in a one-dimensional well is a half-period sine wave which goes to zero at the two edges of the well. The energy of the particle is given by: where is the Planck constant, is the mass of the particle, is the energy state ( corresponds to the ground-state energy), and is the width of the well. Quantum field theory In quantum field theory (QFT), the fabric of "empty" space is visualized as consisting of fields, with the field at every point in space and time being a quantum harmonic oscillator, with neighboring oscillators interacting with each other. According to QFT the universe is made up of matter fields whose quanta are fermions (e.g. electrons and quarks), force fields whose quanta are bosons (i.e. photons and gluons) and a Higgs field whose quantum is the Higgs boson. The matter and force fields have zero-point energy. A related term is zero-point field (ZPF), which is the lowest energy state of a particular field. The vacuum can be viewed not as empty space, but as the combination of all zero-point fields. In QFT the zero-point energy of the vacuum state is called the vacuum energy and the average expectation value of the Hamiltonian is called the vacuum expectation value (also called condensate or simply VEV). The QED vacuum is a part of the vacuum state which specifically deals with quantum electrodynamics (e.g. electromagnetic interactions between photons, electrons and the vacuum) and the QCD vacuum deals with quantum chromodynamics (e.g. color charge interactions between quarks, gluons and the vacuum). Recent experiments advocate the idea that particles themselves can be thought of as excited states of the underlying quantum vacuum, and that all properties of matter are merely vacuum fluctuations arising from interactions with the zero-point field. Each point in space makes a contribution of , resulting in a calculation of infinite zero-point energy in any finite volume; this is one reason renormalization is needed to make sense of quantum field theories. In cosmology, the vacuum energy is one possible explanation for the cosmological constant and the source of dark energy. Scientists are not in agreement about how much energy is contained in the vacuum. Quantum mechanics requires the energy to be large as Paul Dirac claimed it is, like a sea of energy. Other scientists specializing in General Relativity require the energy to be small enough for curvature of space to agree with observed astronomy. The Heisenberg uncertainty principle allows the energy to be as large as needed to promote quantum actions for a brief moment of time, even if the average energy is small enough to satisfy relativity and flat space. To cope with disagreements, the vacuum energy is described as a virtual energy potential of positive and negative energy. In quantum perturbation theory, it is sometimes said that the contribution of one-loop and multi-loop Feynman diagrams to elementary particle propagators are the contribution of vacuum fluctuations, or the zero-point energy to the particle masses. Quantum electrodynamic vacuum The oldest and best known quantized force field is the electromagnetic field. Maxwell's equations have been superseded by quantum electrodynamics (QED). By considering the zero-point energy that arises from QED it is possible to gain a characteristic understanding of zero-point energy that arises not just through electromagnetic interactions but in all quantum field theories. Redefining the zero of energy In the quantum theory of the electromagnetic field, classical wave amplitudes and are replaced by operators and that satisfy: The classical quantity appearing in the classical expression for the energy of a field mode is replaced in quantum theory by the photon number operator . The fact that: implies that quantum theory does not allow states of the radiation field for which the photon number and a field amplitude can be precisely defined, i.e., we cannot have simultaneous eigenstates for and . The reconciliation of wave and particle attributes of the field is accomplished via the association of a probability amplitude with a classical mode pattern. The calculation of field modes is entirely classical problem, while the quantum properties of the field are carried by the mode "amplitudes" and associated with these classical modes. The zero-point energy of the field arises formally from the non-commutativity of and . This is true for any harmonic oscillator: the zero-point energy appears when we write the Hamiltonian: It is often argued that the entire universe is completely bathed in the zero-point electromagnetic field, and as such it can add only some constant amount to expectation values. Physical measurements will therefore reveal only deviations from the vacuum state. Thus the zero-point energy can be dropped from the Hamiltonian by redefining the zero of energy, or by arguing that it is a constant and therefore has no effect on Heisenberg equations of motion. Thus we can choose to declare by fiat that the ground state has zero energy and a field Hamiltonian, for example, can be replaced by: without affecting any physical predictions of the theory. The new Hamiltonian is said to be normally ordered (or Wick ordered) and is denoted by a double-dot symbol. The normally ordered Hamiltonian is denoted , i.e.: In other words, within the normal ordering symbol we can commute and . Since zero-point energy is intimately connected to the non-commutativity of and , the normal ordering procedure eliminates any contribution from the zero-point field. This is especially reasonable in the case of the field Hamiltonian, since the zero-point term merely adds a constant energy which can be eliminated by a simple redefinition for the zero of energy. Moreover, this constant energy in the Hamiltonian obviously commutes with and and so cannot have any effect on the quantum dynamics described by the Heisenberg equations of motion. However, things are not quite that simple. The zero-point energy cannot be eliminated by dropping its energy from the Hamiltonian: When we do this and solve the Heisenberg equation for a field operator, we must include the vacuum field, which is the homogeneous part of the solution for the field operator. In fact we can show that the vacuum field is essential for the preservation of the commutators and the formal consistency of QED. When we calculate the field energy we obtain not only a contribution from particles and forces that may be present but also a contribution from the vacuum field itself i.e. the zero-point field energy. In other words, the zero-point energy reappears even though we may have deleted it from the Hamiltonian. Electromagnetic field in free space From Maxwell's equations, the electromagnetic energy of a "free" field i.e. one with no sources, is described by: We introduce the "mode function" that satisfies the Helmholtz equation: where and assume it is normalized such that: We wish to "quantize" the electromagnetic energy of free space for a multimode field. The field intensity of free space should be independent of position such that should be independent of for each mode of the field. The mode function satisfying these conditions is: where in order to have the transversality condition satisfied for the Coulomb gauge in which we are working. To achieve the desired normalization we pretend space is divided into cubes of volume and impose on the field the periodic boundary condition: or equivalently where can assume any integer value. This allows us to consider the field in any one of the imaginary cubes and to define the mode function: which satisfies the Helmholtz equation, transversality, and the "box normalization": where is chosen to be a unit vector which specifies the polarization of the field mode. The condition means that there are two independent choices of , which we call and where and . Thus we define the mode functions: in terms of which the vector potential becomes: or: where and , are photon annihilation and creation operators for the mode with wave vector and polarization . This gives the vector potential for a plane wave mode of the field. The condition for shows that there are infinitely many such modes. The linearity of Maxwell's equations allows us to write: for the total vector potential in free space. Using the fact that: we find the field Hamiltonian is: This is the Hamiltonian for an infinite number of uncoupled harmonic oscillators. Thus different modes of the field are independent and satisfy the commutation relations: Clearly the least eigenvalue for is: This state describes the zero-point energy of the vacuum. It appears that this sum is divergent – in fact highly divergent, as putting in the density factor shows. The summation becomes approximately the integral: for high values of . It diverges proportional to for large . There are two separate questions to consider. First, is the divergence a real one such that the zero-point energy really is infinite? If we consider the volume is contained by perfectly conducting walls, very high frequencies can only be contained by taking more and more perfect conduction. No actual method of containing the high frequencies is possible. Such modes will not be stationary in our box and thus not countable in the stationary energy content. So from this physical point of view the above sum should only extend to those frequencies which are countable; a cut-off energy is thus eminently reasonable. However, on the scale of a "universe" questions of general relativity must be included. Suppose even the boxes could be reproduced, fit together and closed nicely by curving spacetime. Then exact conditions for running waves may be possible. However the very high frequency quanta will still not be contained. As per John Wheeler's "geons" these will leak out of the system. So again a cut-off is permissible, almost necessary. The question here becomes one of consistency since the very high energy quanta will act as a mass source and start curving the geometry. This leads to the second question. Divergent or not, finite or infinite, is the zero-point energy of any physical significance? The ignoring of the whole zero-point energy is often encouraged for all practical calculations. The reason for this is that energies are not typically defined by an arbitrary data point, but rather changes in data points, so adding or subtracting a constant (even if infinite) should be allowed. However this is not the whole story, in reality energy is not so arbitrarily defined: in general relativity the seat of the curvature of spacetime is the energy content and there the absolute amount of energy has real physical meaning. There is no such thing as an arbitrary additive constant with density of field energy. Energy density curves space, and an increase in energy density produces an increase of curvature. Furthermore, the zero-point energy density has other physical consequences e.g. the Casimir effect, contribution to the Lamb shift, or anomalous magnetic moment of the electron, it is clear it is not just a mathematical constant or artifact that can be cancelled out. Necessity of the vacuum field in QED The vacuum state of the "free" electromagnetic field (that with no sources) is defined as the ground state in which for all modes . The vacuum state, like all stationary states of the field, is an eigenstate of the Hamiltonian but not the electric and magnetic field operators. In the vacuum state, therefore, the electric and magnetic fields do not have definite values. We can imagine them to be fluctuating about their mean value of zero. In a process in which a photon is annihilated (absorbed), we can think of the photon as making a transition into the vacuum state. Similarly, when a photon is created (emitted), it is occasionally useful to imagine that the photon has made a transition out of the vacuum state. An atom, for instance, can be considered to be "dressed" by emission and reabsorption of "virtual photons" from the vacuum. The vacuum state energy described by is infinite. We can make the replacement: the zero-point energy density is: or in other words the spectral energy density of the vacuum field: The zero-point energy density in the frequency range from to is therefore: This can be large even in relatively narrow "low frequency" regions of the spectrum. In the optical region from 400 to 700 nm, for instance, the above equation yields around 220 erg/cm3. We showed in the above section that the zero-point energy can be eliminated from the Hamiltonian by the normal ordering prescription. However, this elimination does not mean that the vacuum field has been rendered unimportant or without physical consequences. To illustrate this point we consider a linear dipole oscillator in the vacuum. The Hamiltonian for the oscillator plus the field with which it interacts is: This has the same form as the corresponding classical Hamiltonian and the Heisenberg equations of motion for the oscillator and the field are formally the same as their classical counterparts. For instance the Heisenberg equations for the coordinate and the canonical momentum of the oscillator are: or: since the rate of change of the vector potential in the frame of the moving charge is given by the convective derivative For nonrelativistic motion we may neglect the magnetic force and replace the expression for by: Above we have made the electric dipole approximation in which the spatial dependence of the field is neglected. The Heisenberg equation for is found similarly from the Hamiltonian to be: in the electric dipole approximation. In deriving these equations for , , and we have used the fact that equal-time particle and field operators commute. This follows from the assumption that particle and field operators commute at some time (say, ) when the matter-field interpretation is presumed to begin, together with the fact that a Heisenberg-picture operator evolves in time as , where is the time evolution operator satisfying Alternatively, we can argue that these operators must commute if we are to obtain the correct equations of motion from the Hamiltonian, just as the corresponding Poisson brackets in classical theory must vanish in order to generate the correct Hamilton equations. The formal solution of the field equation is: and therefore the equation for may be written: where and It can be shown that in the radiation reaction field, if the mass is regarded as the "observed" mass then we can take The total field acting on the dipole has two parts, and . is the free or zero-point field acting on the dipole. It is the homogeneous solution of the Maxwell equation for the field acting on the dipole, i.e., the solution, at the position of the dipole, of the wave equation satisfied by the field in the (source free) vacuum. For this reason is often referred to as the "vacuum field", although it is of course a Heisenberg-picture operator acting on whatever state of the field happens to be appropriate at . is the source field, the field generated by the dipole and acting on the dipole. Using the above equation for we obtain an equation for the Heisenberg-picture operator that is formally the same as the classical equation for a linear dipole oscillator: where . in this instance we have considered a dipole in the vacuum, without any "external" field acting on it. the role of the external field in the above equation is played by the vacuum electric field acting on the dipole. Classically, a dipole in the vacuum is not acted upon by any "external" field: if there are no sources other than the dipole itself, then the only field acting on the dipole is its own radiation reaction field. In quantum theory however there is always an "external" field, namely the source-free or vacuum field . According to our earlier equation for the free field is the only field in existence at as the time at which the interaction between the dipole and the field is "switched on". The state vector of the dipole-field system at is therefore of the form where is the vacuum state of the field and is the initial state of the dipole oscillator. The expectation value of the free field is therefore at all times equal to zero: since . however, the energy density associated with the free field is infinite: The important point of this is that the zero-point field energy does not affect the Heisenberg equation for since it is a c-number or constant (i.e. an ordinary number rather than an operator) and commutes with . We can therefore drop the zero-point field energy from the Hamiltonian, as is usually done. But the zero-point field re-emerges as the homogeneous solution for the field equation. A charged particle in the vacuum will therefore always see a zero-point field of infinite density. This is the origin of one of the infinities of quantum electrodynamics, and it cannot be eliminated by the trivial expedient dropping of the term in the field Hamiltonian. The free field is in fact necessary for the formal consistency of the theory. In particular, it is necessary for the preservation of the commutation relations, which is required by the unitary of time evolution in quantum theory: We can calculate from the formal solution of the operator equation of motion Using the fact that and that equal-time particle and field operators commute, we obtain: For the dipole oscillator under consideration it can be assumed that the radiative damping rate is small compared with the natural oscillation frequency, i.e., . Then the integrand above is sharply peaked at and: the necessity of the vacuum field can also be appreciated by making the small damping approximation in and Without the free field in this equation the operator would be exponentially dampened, and commutators like would approach zero for . With the vacuum field included, however, the commutator is at all times, as required by unitarity, and as we have just shown. A similar result is easily worked out for the case of a free particle instead of a dipole oscillator. What we have here is an example of a "fluctuation-dissipation elation". Generally speaking if a system is coupled to a bath that can take energy from the system in an effectively irreversible way, then the bath must also cause fluctuations. The fluctuations and the dissipation go hand in hand we cannot have one without the other. In the current example the coupling of a dipole oscillator to the electromagnetic field has a dissipative component, in the form of the zero-point (vacuum) field; given the existence of radiation reaction, the vacuum field must also exist in order to preserve the canonical commutation rule and all it entails. The spectral density of the vacuum field is fixed by the form of the radiation reaction field, or vice versa: because the radiation reaction field varies with the third derivative of , the spectral energy density of the vacuum field must be proportional to the third power of in order for to hold. In the case of a dissipative force proportional to , by contrast, the fluctuation force must be proportional to in order to maintain the canonical commutation relation. This relation between the form of the dissipation and the spectral density of the fluctuation is the essence of the fluctuation-dissipation theorem. The fact that the canonical commutation relation for a harmonic oscillator coupled to the vacuum field is preserved implies that the zero-point energy of the oscillator is preserved. it is easy to show that after a few damping times the zero-point motion of the oscillator is in fact sustained by the driving zero-point field. Quantum chromodynamic vacuum The QCD vacuum is the vacuum state of quantum chromodynamics (QCD). It is an example of a non-perturbative vacuum state, characterized by a non-vanishing condensates such as the gluon condensate and the quark condensate in the complete theory which includes quarks. The presence of these condensates characterizes the confined phase of quark matter. In technical terms, gluons are vector gauge bosons that mediate strong interactions of quarks in quantum chromodynamics (QCD). Gluons themselves carry the color charge of the strong interaction. This is unlike the photon, which mediates the electromagnetic interaction but lacks an electric charge. Gluons therefore participate in the strong interaction in addition to mediating it, making QCD significantly harder to analyze than QED (quantum electrodynamics) as it deals with nonlinear equations to characterize such interactions. Higgs field The Standard Model hypothesises a field called the Higgs field (symbol: ), which has the unusual property of a non-zero amplitude in its ground state (zero-point) energy after renormalization; i.e., a non-zero vacuum expectation value. It can have this effect because of its unusual "Mexican hat" shaped potential whose lowest "point" is not at its "centre". Below a certain extremely high energy level the existence of this non-zero vacuum expectation spontaneously breaks electroweak gauge symmetry which in turn gives rise to the Higgs mechanism and triggers the acquisition of mass by those particles interacting with the field. The Higgs mechanism occurs whenever a charged field has a vacuum expectation value. This effect occurs because scalar field components of the Higgs field are "absorbed" by the massive bosons as degrees of freedom, and couple to the fermions via Yukawa coupling, thereby producing the expected mass terms. The expectation value of in the ground state (the vacuum expectation value or VEV) is then , where . The measured value of this parameter is approximately . It has units of mass, and is the only free parameter of the Standard Model that is not a dimensionless number. The Higgs mechanism is a type of superconductivity which occurs in the vacuum. It occurs when all of space is filled with a sea of particles which are charged and thus the field has a nonzero vacuum expectation value. Interaction with the vacuum energy filling the space prevents certain forces from propagating over long distances (as it does in a superconducting medium; e.g., in the Ginzburg–Landau theory). Experimental observations Zero-point energy has many observed physical consequences. It is important to note that zero-point energy is not merely an artifact of mathematical formalism that can, for instance, be dropped from a Hamiltonian by redefining the zero of energy, or by arguing that it is a constant and therefore has no effect on Heisenberg equations of motion without latter consequence. Indeed, such treatment could create a problem at a deeper, as of yet undiscovered, theory. For instance, in general relativity the zero of energy (i.e. the energy density of the vacuum) contributes to a cosmological constant of the type introduced by Einstein in order to obtain static solutions to his field equations. The zero-point energy density of the vacuum, due to all quantum fields, is extremely large, even when we cut off the largest allowable frequencies based on plausible physical arguments. It implies a cosmological constant larger than the limits imposed by observation by about 120 orders of magnitude. This "cosmological constant problem" remains one of the greatest unsolved mysteries of physics. Casimir effect A phenomenon that is commonly presented as evidence for the existence of zero-point energy in vacuum is the Casimir effect, proposed in 1948 by Dutch physicist Hendrik Casimir, who considered the quantized electromagnetic field between a pair of grounded, neutral metal plates. The vacuum energy contains contributions from all wavelengths, except those excluded by the spacing between plates. As the plates draw together, more wavelengths are excluded and the vacuum energy decreases. The decrease in energy means there must be a force doing work on the plates as they move. Early experimental tests from the 1950s onwards gave positive results showing the force was real, but other external factors could not be ruled out as the primary cause, with the range of experimental error sometimes being nearly 100%. That changed in 1997 with Lamoreaux conclusively showing that the Casimir force was real. Results have been repeatedly replicated since then. In 2009, Munday et al. published experimental proof that (as predicted in 1961) the Casimir force could also be repulsive as well as being attractive. Repulsive Casimir forces could allow quantum levitation of objects in a fluid and lead to a new class of switchable nanoscale devices with ultra-low static friction. An interesting hypothetical side effect of the Casimir effect is the Scharnhorst effect, a hypothetical phenomenon in which light signals travel slightly faster than between two closely spaced conducting plates. Lamb shift The quantum fluctuations of the electromagnetic field have important physical consequences. In addition to the Casimir effect, they also lead to a splitting between the two energy levels and (in term symbol notation) of the hydrogen atom which was not predicted by the Dirac equation, according to which these states should have the same energy. Charged particles can interact with the fluctuations of the quantized vacuum field, leading to slight shifts in energy; this effect is called the Lamb shift. The shift of about is roughly of the difference between the energies of the 1s and 2s levels, and amounts to 1,058 MHz in frequency units. A small part of this shift (27 MHz ≈ 3%) arises not from fluctuations of the electromagnetic field, but from fluctuations of the electron–positron field. The creation of (virtual) electron–positron pairs has the effect of screening the Coulomb field and acts as a vacuum dielectric constant. This effect is much more important in muonic atoms. Fine-structure constant Taking (the Planck constant divided by ), (the speed of light), and (the electromagnetic coupling constant i.e. a measure of the strength of the electromagnetic force (where is the absolute value of the electronic charge and is the vacuum permittivity)) we can form a dimensionless quantity called the fine-structure constant: The fine-structure constant is the coupling constant of quantum electrodynamics (QED) determining the strength of the interaction between electrons and photons. It turns out that the fine-structure constant is not really a constant at all owing to the zero-point energy fluctuations of the electron-positron field. The quantum fluctuations caused by zero-point energy have the effect of screening electric charges: owing to (virtual) electron-positron pair production, the charge of the particle measured far from the particle is far smaller than the charge measured when close to it. The Heisenberg inequality where , and , are the standard deviations of position and momentum states that: It means that a short distance implies large momentum and therefore high energy i.e. particles of high energy must be used to explore short distances. QED concludes that the fine-structure constant is an increasing function of energy. It has been shown that at energies of the order of the Z0 boson rest energy, 90 GeV, that: rather than the low-energy . The renormalization procedure of eliminating zero-point energy infinities allows the choice of an arbitrary energy (or distance) scale for defining . All in all, depends on the energy scale characteristic of the process under study, and also on details of the renormalization procedure. The energy dependence of has been observed for several years now in precision experiment in high-energy physics. Vacuum birefringence In the presence of strong electrostatic fields it is predicted that virtual particles become separated from the vacuum state and form real matter. The fact that electromagnetic radiation can be transformed into matter and vice versa leads to fundamentally new features in quantum electrodynamics. One of the most important consequences is that, even in the vacuum, the Maxwell equations have to be exchanged by more complicated formulas. In general, it will be not possible to separate processes in the vacuum from the processes involving matter since electromagnetic fields can create matter if the field fluctuations are strong enough. This leads to highly complex nonlinear interaction – gravity will have an effect on the light at the same time the light has an effect on gravity. These effects were first predicted by Werner Heisenberg and Hans Heinrich Euler in 1936 and independently the same year by Victor Weisskopf who stated: "The physical properties of the vacuum originate in the "zero-point energy" of matter, which also depends on absent particles through the external field strengths and therefore contributes an additional term to the purely Maxwellian field energy". Thus strong magnetic fields vary the energy contained in the vacuum. The scale above which the electromagnetic field is expected to become nonlinear is known as the Schwinger limit. At this point the vacuum has all the properties of a birefringent medium, thus in principle a rotation of the polarization frame (the Faraday effect) can be observed in empty space. Both Einstein's theory of special and general relativity state that light should pass freely through a vacuum without being altered, a principle known as Lorentz invariance. Yet, in theory, large nonlinear self-interaction of light due to quantum fluctuations should lead to this principle being measurably violated if the interactions are strong enough. Nearly all theories of quantum gravity predict that Lorentz invariance is not an exact symmetry of nature. It is predicted the speed at which light travels through the vacuum depends on its direction, polarization and the local strength of the magnetic field. There have been a number of inconclusive results which claim to show evidence of a Lorentz violation by finding a rotation of the polarization plane of light coming from distant galaxies. The first concrete evidence for vacuum birefringence was published in 2017 when a team of astronomers looked at the light coming from the star RX J1856.5-3754, the closest discovered neutron star to Earth. Roberto Mignani at the National Institute for Astrophysics in Milan who led the team of astronomers has commented that "When Einstein came up with the theory of general relativity 100 years ago, he had no idea that it would be used for navigational systems. The consequences of this discovery probably will also have to be realised on a longer timescale." The team found that visible light from the star had undergone linear polarisation of around 16%. If the birefringence had been caused by light passing through interstellar gas or plasma, the effect should have been no more than 1%. Definitive proof would require repeating the observation at other wavelengths and on other neutron stars. At X-ray wavelengths the polarization from the quantum fluctuations should be near 100%. Although no telescope currently exists that can make such measurements, there are several proposed X-ray telescopes that may soon be able to verify the result conclusively such as China's Hard X-ray Modulation Telescope (HXMT) and NASA's Imaging X-ray Polarimetry Explorer (IXPE). Speculated involvement in other phenomena Dark energy In the late 1990s it was discovered that very distant supernovae were dimmer than expected suggesting that the universe's expansion was accelerating rather than slowing down. This revived discussion that Einstein's cosmological constant, long disregarded by physicists as being equal to zero, was in fact some small positive value. This would indicate empty space exerted some form of negative pressure or energy. There is no natural candidate for what might cause what has been called dark energy but the current best guess is that it is the zero-point energy of the vacuum, but this guess is known to be off by 120 orders of magnitude. The European Space Agency's Euclid telescope, launched on 1 July 2023, will map galaxies up to 10 billion light years away. By seeing how dark energy influences their arrangement and shape, the mission will allow scientists to see if the strength of dark energy has changed. If dark energy is found to vary throughout time it would indicate it is due to quintessence, where observed acceleration is due to the energy of a scalar field, rather than the cosmological constant. No evidence of quintessence is yet available, but it has not been ruled out either. It generally predicts a slightly slower acceleration of the expansion of the universe than the cosmological constant. Some scientists think that the best evidence for quintessence would come from violations of Einstein's equivalence principle and variation of the fundamental constants in space or time. Scalar fields are predicted by the Standard Model of particle physics and string theory, but an analogous problem to the cosmological constant problem (or the problem of constructing models of cosmological inflation) occurs: renormalization theory predicts that scalar fields should acquire large masses again due to zero-point energy. Cosmic inflation Cosmic inflation is phase of accelerated cosmic expansion just after the Big Bang. It explains the origin of the large-scale structure of the cosmos. It is believed quantum vacuum fluctuations caused by zero-point energy arising in the microscopic inflationary period, later became magnified to a cosmic size, becoming the gravitational seeds for galaxies and structure in the Universe (see galaxy formation and evolution and structure formation). Many physicists also believe that inflation explains why the Universe appears to be the same in all directions (isotropic), why the cosmic microwave background radiation is distributed evenly, why the Universe is flat, and why no magnetic monopoles have been observed. The mechanism for inflation is unclear, it is similar in effect to dark energy but is a far more energetic and short lived process. As with dark energy the best explanation is some form of vacuum energy arising from quantum fluctuations. It may be that inflation caused baryogenesis, the hypothetical physical processes that produced an asymmetry (imbalance) between baryons and antibaryons produced in the very early universe, but this is far from certain. Cosmology Paul S. Wesson examined the cosmological implications of assuming that zero-point energy is real. Among numerous difficulties, general relativity requires that such energy not gravitate, so it cannot be similar to electromagnetic radiation. Alternative theories There has been a long debate over the question of whether zero-point fluctuations of quantized vacuum fields are "real" i.e. do they have physical effects that cannot be interpreted by an equally valid alternative theory? Schwinger, in particular, attempted to formulate QED without reference to zero-point fluctuations via his "source theory". From such an approach it is possible to derive the Casimir Effect without reference to a fluctuating field. Such a derivation was first given by Schwinger (1975) for a scalar field, and then generalized to the electromagnetic case by Schwinger, DeRaad, and Milton (1978). in which they state "the vacuum is regarded as truly a state with all physical properties equal to zero". Jaffe (2005) has highlighted a similar approach in deriving the Casimir effect stating "the concept of zero-point fluctuations is a heuristic and calculational aid in the description of the Casimir effect, but not a necessity in QED." Milonni has shown the necessity of the vacuum field for the formal consistency of QED. Modern physics does not know any better way to construct gauge-invariant, renormalizable theories than with zero-point energy and they would seem to be a necessity for any attempt at a unified theory. Nevertheless, as pointed out by Jaffe, "no known phenomenon, including the Casimir effect, demonstrates that zero point energies are “real”" Chaotic and emergent phenomena The mathematical models used in classical electromagnetism, quantum electrodynamics (QED) and the Standard Model all view the electromagnetic vacuum as a linear system with no overall observable consequence. For example, in the case of the Casimir effect, Lamb shift, and so on these phenomena can be explained by alternative mechanisms other than action of the vacuum by arbitrary changes to the normal ordering of field operators. See the alternative theories section. This is a consequence of viewing electromagnetism as a U(1) gauge theory, which topologically does not allow the complex interaction of a field with and on itself. In higher symmetry groups and in reality, the vacuum is not a calm, randomly fluctuating, largely immaterial and passive substance, but at times can be viewed as a turbulent virtual plasma that can have complex vortices (i.e. solitons vis-à-vis particles), entangled states and a rich nonlinear structure. There are many observed nonlinear physical electromagnetic phenomena such as Aharonov–Bohm (AB) and Altshuler–Aronov–Spivak (AAS) effects, Berry, Aharonov–Anandan, Pancharatnam and Chiao–Wu phase rotation effects, Josephson effect, Quantum Hall effect, the De Haas–Van Alphen effect, the Sagnac effect and many other physically observable phenomena which would indicate that the electromagnetic potential field has real physical meaning rather than being a mathematical artifact and therefore an all encompassing theory would not confine electromagnetism as a local force as is currently done, but as a SU(2) gauge theory or higher geometry. Higher symmetries allow for nonlinear, aperiodic behaviour which manifest as a variety of complex non-equilibrium phenomena that do not arise in the linearised U(1) theory, such as multiple stable states, symmetry breaking, chaos and emergence. What are called Maxwell's equations today, are in fact a simplified version of the original equations reformulated by Heaviside, FitzGerald, Lodge and Hertz. The original equations used Hamilton's more expressive quaternion notation, a kind of Clifford algebra, which fully subsumes the standard Maxwell vectorial equations largely used today. In the late 1880s there was a debate over the relative merits of vector analysis and quaternions. According to Heaviside the electromagnetic potential field was purely metaphysical, an arbitrary mathematical fiction, that needed to be "murdered". It was concluded that there was no need for the greater physical insights provided by the quaternions if the theory was purely local in nature. Local vector analysis has become the dominant way of using Maxwell's equations ever since. However, this strictly vectorial approach has led to a restrictive topological understanding in some areas of electromagnetism, for example, a full understanding of the energy transfer dynamics in Tesla's oscillator-shuttle-circuit can only be achieved in quaternionic algebra or higher SU(2) symmetries. It has often been argued that quaternions are not compatible with special relativity, but multiple papers have shown ways of incorporating relativity. A good example of nonlinear electromagnetics is in high energy dense plasmas, where vortical phenomena occur which seemingly violate the second law of thermodynamics by increasing the energy gradient within the electromagnetic field and violate Maxwell's laws by creating ion currents which capture and concentrate their own and surrounding magnetic fields. In particular Lorentz force law, which elaborates Maxwell's equations is violated by these force free vortices. These apparent violations are due to the fact that the traditional conservation laws in classical and quantum electrodynamics (QED) only display linear U(1) symmetry (in particular, by the extended Noether theorem, conservation laws such as the laws of thermodynamics need not always apply to dissipative systems, which are expressed in gauges of higher symmetry). The second law of thermodynamics states that in a closed linear system entropy flow can only be positive (or exactly zero at the end of a cycle). However, negative entropy (i.e. increased order, structure or self-organisation) can spontaneously appear in an open nonlinear thermodynamic system that is far from equilibrium, so long as this emergent order accelerates the overall flow of entropy in the total system. The 1977 Nobel Prize in Chemistry was awarded to thermodynamicist Ilya Prigogine for his theory of dissipative systems that described this notion. Prigogine described the principle as "order through fluctuations" or "order out of chaos". It has been argued by some that all emergent order in the universe from galaxies, solar systems, planets, weather, complex chemistry, evolutionary biology to even consciousness, technology and civilizations are themselves examples of thermodynamic dissipative systems; nature having naturally selected these structures to accelerate entropy flow within the universe to an ever-increasing degree. For example, it has been estimated that human body is 10,000 times more effective at dissipating energy per unit of mass than the sun. One may query what this has to do with zero-point energy. Given the complex and adaptive behaviour that arises from nonlinear systems considerable attention in recent years has gone into studying a new class of phase transitions which occur at absolute zero temperature. These are quantum phase transitions which are driven by EM field fluctuations as a consequence of zero-point energy. A good example of a spontaneous phase transition that are attributed to zero-point fluctuations can be found in superconductors. Superconductivity is one of the best known empirically quantified macroscopic electromagnetic phenomena whose basis is recognised to be quantum mechanical in origin. The behaviour of the electric and magnetic fields under superconductivity is governed by the London equations. However, it has been questioned in a series of journal articles whether the quantum mechanically canonised London equations can be given a purely classical derivation. Bostick, for instance, has claimed to show that the London equations do indeed have a classical origin that applies to superconductors and to some collisionless plasmas as well. In particular it has been asserted that the Beltrami vortices in the plasma focus display the same paired flux-tube morphology as Type II superconductors. Others have also pointed out this connection, Fröhlich has shown that the hydrodynamic equations of compressible fluids, together with the London equations, lead to a macroscopic parameter ( = electric charge density / mass density), without involving either quantum phase factors or the Planck constant. In essence, it has been asserted that Beltrami plasma vortex structures are able to at least simulate the morphology of Type I and Type II superconductors. This occurs because the "organised" dissipative energy of the vortex configuration comprising the ions and electrons far exceeds the "disorganised" dissipative random thermal energy. The transition from disorganised fluctuations to organised helical structures is a phase transition involving a change in the condensate's energy (i.e. the ground state or zero-point energy) but without any associated rise in temperature. This is an example of zero-point energy having multiple stable states (see Quantum phase transition, Quantum critical point, Topological degeneracy, Topological order) and where the overall system structure is independent of a reductionist or deterministic view, that "classical" macroscopic order can also causally affect quantum phenomena. Furthermore, the pair production of Beltrami vortices has been compared to the morphology of pair production of virtual particles in the vacuum. The idea that the vacuum energy can have multiple stable energy states is a leading hypothesis for the cause of cosmic inflation. In fact, it has been argued that these early vacuum fluctuations led to the expansion of the universe and in turn have guaranteed the non-equilibrium conditions necessary to drive order from chaos, as without such expansion the universe would have reached thermal equilibrium and no complexity could have existed. With the continued accelerated expansion of the universe, the cosmos generates an energy gradient that increases the "free energy" (i.e. the available, usable or potential energy for useful work) which the universe is able to use to create ever more complex forms of order. The only reason Earth's environment does not decay into an equilibrium state is that it receives a daily dose of sunshine and that, in turn, is due to the sun "polluting" interstellar space with entropy. The sun's fusion power is only possible due to the gravitational disequilibrium of matter that arose from cosmic expansion. In this essence, the vacuum energy can be viewed as the key cause of the structure throughout the universe. That humanity might alter the morphology of the vacuum energy to create an energy gradient for useful work is the subject of much controversy. Purported applications Physicists overwhelmingly reject any possibility that the zero-point energy field can be exploited to obtain useful energy (work) or uncompensated momentum; such efforts are seen as tantamount to perpetual motion machines. Nevertheless, the allure of free energy has motivated such research, usually falling in the category of fringe science. As long ago as 1889 (before quantum theory or discovery of the zero point energy) Nikola Tesla proposed that useful energy could be obtained from free space, or what was assumed at that time to be an all-pervasive aether. Others have since claimed to exploit zero-point or vacuum energy with a large amount of pseudoscientific literature causing ridicule around the subject. Despite rejection by the scientific community, harnessing zero-point energy remains an interest of research, particularly in the US where it has attracted the attention of major aerospace/defence contractors and the U.S. Department of Defense as well as in China, Germany, Russia and Brazil. Casimir batteries and engines A common assumption is that the Casimir force is of little practical use; the argument is made that the only way to actually gain energy from the two plates is to allow them to come together (getting them apart again would then require more energy), and therefore it is a one-use-only tiny force in nature. In 1984 Robert Forward published work showing how a "vacuum-fluctuation battery" could be constructed; the battery can be recharged by making the electrical forces slightly stronger than the Casimir force to reexpand the plates. In 1999, Pinto, a former scientist at NASA's Jet Propulsion Laboratory at Caltech in Pasadena, published in Physical Review his thought experiment (Gedankenexperiment) for a "Casimir engine". The paper showed that continuous positive net exchange of energy from the Casimir effect was possible, even stating in the abstract "In the event of no other alternative explanations, one should conclude that major technological advances in the area of endless, by-product free-energy production could be achieved." Garret Moddel at University of Colorado has highlighted that he believes such devices hinge on the assumption that the Casimir force is a nonconservative force, he argues that there is sufficient evidence (e.g. analysis by Scandurra (2001)) to say that the Casimir effect is a conservative force and therefore even though such an engine can exploit the Casimir force for useful work it cannot produce more output energy than has been input into the system. In 2008, DARPA solicited research proposals in the area of Casimir Effect Enhancement (CEE). The goal of the program is to develop new methods to control and manipulate attractive and repulsive forces at surfaces based on engineering of the Casimir force. A 2008 patent by Haisch and Moddel details a device that is able to extract power from zero-point fluctuations using a gas that circulates through a Casimir cavity. A published test of this concept by Moddel was performed in 2012 and seemed to give excess energy that could not be attributed to another source. However it has not been conclusively shown to be from zero-point energy and the theory requires further investigation. Single heat baths In 1951 Callen and Welton proved the quantum fluctuation-dissipation theorem (FDT) which was originally formulated in classical form by Nyquist (1928) as an explanation for observed Johnson noise in electric circuits. Fluctuation-dissipation theorem showed that when something dissipates energy, in an effectively irreversible way, a connected heat bath must also fluctuate. The fluctuations and the dissipation go hand in hand; it is impossible to have one without the other. The implication of FDT being that the vacuum could be treated as a heat bath coupled to a dissipative force and as such energy could, in part, be extracted from the vacuum for potentially useful work. Such a theory has met with resistance: Macdonald (1962) and Harris (1971) claimed that extracting power from the zero-point energy to be impossible, so FDT could not be true. Grau and Kleen (1982) and Kleen (1986), argued that the Johnson noise of a resistor connected to an antenna must satisfy Planck's thermal radiation formula, thus the noise must be zero at zero temperature and FDT must be invalid. Kiss (1988) pointed out that the existence of the zero-point term may indicate that there is a renormalization problem—i.e., a mathematical artifact—producing an unphysical term that is not actually present in measurements (in analogy with renormalization problems of ground states in quantum electrodynamics). Later, Abbott et al. (1996) arrived at a different but unclear conclusion that "zero-point energy is infinite thus it should be renormalized but not the 'zero-point fluctuations'". Despite such criticism, FDT has been shown to be true experimentally under certain quantum, non-classical conditions. Zero-point fluctuations can, and do, contribute towards systems which dissipate energy. A paper by Armen Allahverdyan and Theo Nieuwenhuizen in 2000 showed the feasibility of extracting zero-point energy for useful work from a single bath, without contradicting the laws of thermodynamics, by exploiting certain quantum mechanical properties. There have been a growing number of papers showing that in some instances the classical laws of thermodynamics, such as limits on the Carnot efficiency, can be violated by exploiting negative entropy of quantum fluctuations. Despite efforts to reconcile quantum mechanics and thermodynamics over the years, their compatibility is still an open fundamental problem. The full extent that quantum properties can alter classical thermodynamic bounds is unknown Space travel and gravitational shielding The use of zero-point energy for space travel is speculative and does not form part of the mainstream scientific consensus. A complete quantum theory of gravitation (that would deal with the role of quantum phenomena like zero-point energy) does not yet exist. Speculative papers explaining a relationship between zero-point energy and gravitational shielding effects have been proposed, but the interaction (if any) is not yet fully understood. According to the general theory of relativity, rotating matter can generate a new force of nature, known as the gravitomagnetic interaction, whose intensity is proportional to the rate of spin. In certain conditions the gravitomagnetic field can be repulsive. In neutron stars for example, it can produce a gravitational analogue of the Meissner effect, but the force produced in such an example is theorized to be exceedingly weak. In 1963 Robert Forward, a physicist and aerospace engineer at Hughes Research Laboratories, published a paper showing how within the framework of general relativity "anti-gravitational" effects might be achieved. Since all atoms have spin, gravitational permeability may be able to differ from material to material. A strong toroidal gravitational field that acts against the force of gravity could be generated by materials that have nonlinear properties that enhance time-varying gravitational fields. Such an effect would be analogous to the nonlinear electromagnetic permeability of iron, making it an effective core (i.e. the doughnut of iron) in a transformer, whose properties are dependent on magnetic permeability. In 1966 Dewitt was first to identify the significance of gravitational effects in superconductors. Dewitt demonstrated that a magnetic-type gravitational field must result in the presence of fluxoid quantization. In 1983, Dewitt's work was substantially expanded by Ross. From 1971 to 1974 Henry William Wallace, a scientist at GE Aerospace was issued with three patents. Wallace used Dewitt's theory to develop an experimental apparatus for generating and detecting a secondary gravitational field, which he named the kinemassic field (now better known as the gravitomagnetic field). In his three patents, Wallace describes three different methods used for detection of the gravitomagnetic field – change in the motion of a body on a pivot, detection of a transverse voltage in a semiconductor crystal, and a change in the specific heat of a crystal material having spin-aligned nuclei. There are no publicly available independent tests verifying Wallace's devices. Such an effect if any would be small. Referring to Wallace's patents, a New Scientist article in 1980 stated "Although the Wallace patents were initially ignored as cranky, observers believe that his invention is now under serious but secret investigation by the military authorities in the USA. The military may now regret that the patents have already been granted and so are available for anyone to read." A further reference to Wallace's patents occur in an electric propulsion study prepared for the Astronautics Laboratory at Edwards Air Force Base which states: "The patents are written in a very believable style which include part numbers, sources for some components, and diagrams of data. Attempts were made to contact Wallace using patent addresses and other sources but he was not located nor is there a trace of what became of his work. The concept can be somewhat justified on general relativistic grounds since rotating frames of time varying fields are expected to emit gravitational waves." In 1986 the U.S. Air Force's then Rocket Propulsion Laboratory (RPL) at Edwards Air Force Base solicited "Non Conventional Propulsion Concepts" under a small business research and innovation program. One of the six areas of interest was "Esoteric energy sources for propulsion, including the quantum dynamic energy of vacuum space..." In the same year BAE Systems launched "Project Greenglow" to provide a "focus for research into novel propulsion systems and the means to power them". In 1988 Kip Thorne et al. published work showing how traversable wormholes can exist in spacetime only if they are threaded by quantum fields generated by some form of exotic matter that has negative energy. In 1993 Scharnhorst and Barton showed that the speed of a photon will be increased if it travels between two Casimir plates, an example of negative energy. In the most general sense, the exotic matter needed to create wormholes would share the repulsive properties of the inflationary energy, dark energy or zero-point radiation of the vacuum. Building on the work of Thorne, in 1994 Miguel Alcubierre proposed a method for changing the geometry of space by creating a wave that would cause the fabric of space ahead of a spacecraft to contract and the space behind it to expand (see Alcubierre drive). The ship would then ride this wave inside a region of flat space, known as a warp bubble and would not move within this bubble but instead be carried along as the region itself moves due to the actions of the drive. In 1992 Evgeny Podkletnov published a heavily debated journal article claiming a specific type of rotating superconductor could shield gravitational force. Independently of this, from 1991 to 1993 Ning Li and Douglas Torr published a number of articles about gravitational effects in superconductors. One finding they derived is the source of gravitomagnetic flux in a type II superconductor material is due to spin alignment of the lattice ions. Quoting from their third paper: "It is shown that the coherent alignment of lattice ion spins will generate a detectable gravitomagnetic field, and in the presence of a time-dependent applied magnetic vector potential field, a detectable gravitoelectric field." The claimed size of the generated force has been disputed by some but defended by others. In 1997 Li published a paper attempting to replicate Podkletnov's results and showed the effect was very small, if it existed at all. Li is reported to have left the University of Alabama in 1999 to found the company AC Gravity LLC. AC Gravity was awarded a U.S. Department of Defense grant for $448,970 in 2001 to continue anti-gravity research. The grant period ended in 2002 but no results from this research were made public. In 2002 Phantom Works, Boeing's advanced research and development facility in Seattle, approached Evgeny Podkletnov directly. Phantom Works was blocked by Russian technology transfer controls. At this time Lieutenant General George Muellner, the outgoing head of the Boeing Phantom Works, confirmed that attempts by Boeing to work with Podkletnov had been blocked by Russian government, lso commenting that "The physical principles – and Podkletnov's device is not the only one – appear to be valid... There is basic science there. They're not breaking the laws of physics. The issue is whether the science can be engineered into something workable" Froning and Roach (2002) put forward a paper that builds on the work of Puthoff, Haisch and Alcubierre. They used fluid dynamic simulations to model the interaction of a vehicle (like that proposed by Alcubierre) with the zero-point field. Vacuum field perturbations are simulated by fluid field perturbations and the aerodynamic resistance of viscous drag exerted on the interior of the vehicle is compared to the Lorentz force exerted by the zero-point field (a Casimir-like force is exerted on the exterior by unbalanced zero-point radiation pressures). They find that the optimized negative energy required for an Alcubierre drive is where it is a saucer-shaped vehicle with toroidal electromagnetic fields. The EM fields distort the vacuum field perturbations surrounding the craft sufficiently to affect the permeability and permittivity of space. In 2009, Giorgio Fontana and Bernd Binder presented a new method to potentially extract the Zero-point energy of the electromagnetic field and nuclear forces in the form of gravitational waves. In the spheron model of the nucleus, proposed by the two times Nobel laureate Linus Pauling, dineutrons are among the components of this structure. Similarly to a dumbbell put in a suitable rotational state, but with nuclear mass density, dineutrons are nearly ideal sources of gravitational waves at X-ray and gamma-ray frequencies. The dynamical interplay, mediated by nuclear forces, between the electrically neutral dineutrons and the electrically charged core nucleus is the fundamental mechanism by which nuclear vibrations can be converted to a rotational state of dineutrons with emission of gravitational waves. Gravity and gravitational waves are well described by General Relativity, that is not a quantum theory, this implies that there is no Zero-point energy for gravity in this theory, therefore dineutrons will emit gravitational waves like any other known source of gravitational waves. In Fontana and Binder paper, nuclear species with dynamical instabilites, related to the Zero-point energy of the electromagnetic field and nuclear forces, and possessing dineutrons, will emit gravitational waves. In experimental physics this approach is still unexplored. In 2014 NASA's Eagleworks Laboratories announced that they had successfully validated the use of a Quantum Vacuum Plasma Thruster which makes use of the Casimir effect for propulsion. In 2016 a scientific paper by the team of NASA scientists passed peer review for the first time. The paper suggests that the zero-point field acts as pilot-wave and that the thrust may be due to particles pushing off the quantum vacuum. While peer review doesn't guarantee that a finding or observation is valid, it does indicate that independent scientists looked over the experimental setup, results, and interpretation and that they could not find any obvious errors in the methodology and that they found the results reasonable. In the paper, the authors identify and discuss nine potential sources of experimental errors, including rogue air currents, leaky electromagnetic radiation, and magnetic interactions. Not all of them could be completely ruled out, and further peer-reviewed experimentation is needed in order to rule these potential errors out. Zero-point energy in fiction The concept of Zero-point energy used as an energy source has been an element used in science fiction and related media.	Physical sciences	Quantum mechanics	Physics
84521	https://en.wikipedia.org/wiki/Tethys%20%28moon%29	Tethys (moon)	Tethys (), or Saturn III, is the fifth-largest moon of Saturn, measuring about across. It was discovered by Giovanni Domenico Cassini in 1684, and is named after the titan Tethys of Greek mythology. Tethys has a low density of 0.98 g/cm3, the lowest of all the major moons in the solar system, indicating that it is made of water ice with just a small fraction of rock. This was confirmed by the spectroscopy of its surface, which identified water ice as the dominant surface material. A further, smaller amount of an unidentified dark material is present as well. The surface of Tethys is very bright, the second-brightest of the moons of Saturn after Enceladus, and neutral in color. Tethys is heavily cratered and cut by a number of large faults and trench-like graben. The largest impact crater, Odysseus, is about 400 km in diameter, whereas the largest graben, Ithaca Chasma, is about 100 km wide and more than 2,000 km long; the two surface features may be related. A small part of the surface is covered by smooth plains that may be cryovolcanic in origin. Like the other regular moons of Saturn, Tethys formed from the Saturnian sub-nebula—a disk of gas and dust that surrounded Saturn soon after its formation. Tethys has been approached and observed by several space probes, including Pioneer 11 (1979), Voyager 1 (1980) and Voyager 2 (1981), with Cassini-Huygens observing the moon the most, and in greatest detail, during its extensive mission to the Saturnian system (2004-2017). Discovery and naming Tethys was discovered by Giovanni Domenico Cassini in 1684 together with Dione, another moon of Saturn. He had also discovered two moons, Rhea and Iapetus earlier, in 1671–72. Cassini observed all of these moons using a large aerial telescope he set up on the grounds of the Paris Observatory. Cassini named the four new moons as Sidera Lodoicea ("the stars of Louis") to honour king Louis XIV of France. By the end of the seventeenth century, astronomers fell into the habit of referring to them and Titan as Saturn I through Saturn V (Tethys, Dione, Rhea, Titan, Iapetus). Once Mimas and Enceladus were discovered in 1789 by William Herschel, the numbering scheme was extended to Saturn VII by bumping the older five moons up two slots. The discovery of Hyperion in 1848 changed the numbers one last time, bumping Iapetus up to Saturn VIII. Henceforth, the numbering scheme would remain fixed. The modern names of all seven satellites of Saturn come from John Herschel (son of William Herschel, discoverer of Mimas and Enceladus). In his 1847 publication Results of Astronomical Observations made at the Cape of Good Hope, he suggested the names of the Titans, sisters and brothers of Kronos (the Greek analogue of Saturn), be used. Tethys is named after the titaness Tethys. It is also designated Saturn III or S III Tethys. The name Tethys has two customary pronunciations, with either a 'long' or a 'short' e: or . The conventional adjectival form of the name is Tethyan, again with either a long or a short e. Orbit Tethys orbits Saturn at a distance of about 295,000 km (about 4.4 Saturn's radii) from the center of the planet. Its orbital eccentricity is negligible, and its orbital inclination is about 1°. Tethys is locked in an inclination resonance with Mimas; however, due to the low gravity of the respective bodies, this interaction does not cause any noticeable orbital eccentricity or tidal heating. The Tethyan orbit lies deep inside the magnetosphere of Saturn, so the plasma co-rotating with the planet strikes the trailing hemisphere of the moon. Tethys is also subject to constant bombardment by the energetic particles (electrons and ions) present in the magnetosphere. Trojans Tethys has two co-orbital moons, Telesto and Calypso, orbiting near Tethys's Lagrange points (60° ahead) and (60° behind) respectively. Physical characteristics Tethys is the 16th-largest moon in the Solar System, with a radius of 531 km. Its mass is about (0.000103 Earth mass), which is less than 1% of the Moon. Despite its relatively small mass, it is more massive than all known moons smaller than itself combined. The density of Tethys is 0.98 g/cm3, indicating that it is composed almost entirely of water-ice. It is also the fifth-largest of Saturn's moons. It is not known whether Tethys is differentiated into a rocky core and ice mantle. However, if it is differentiated, the radius of the core does not exceed 145 km, and its mass is below 6% of the total mass. Due to the action of tidal and rotational forces, Tethys has the shape of triaxial ellipsoid. The dimensions of this ellipsoid are consistent with it having a homogeneous interior. The existence of a subsurface ocean—a layer of liquid salt water in the interior of Tethys—is considered unlikely. The surface of Tethys is one of the most reflective (at visual wavelengths) in the Solar System, with a visual albedo of 1.229. This very high albedo is the result of the sandblasting of particles from Saturn's E-ring, a faint ring composed of small, water-ice particles generated by Enceladus's south polar geysers. The radar albedo of the Tethyan surface is also very high. The leading hemisphere of Tethys is brighter by 10–15% than the trailing one. The high albedo indicates that the surface of Tethys is composed of almost pure water ice with only a small amount of darker materials. The visible spectrum of Tethys is flat and featureless, whereas in the near-infrared strong water ice absorption bands at 1.25, 1.5, 2.0 and 3.0 μm wavelengths are visible. No compound other than crystalline water ice has been unambiguously identified on Tethys. (Possible constituents include organics, ammonia and carbon dioxide.) The dark material in the ice has the same spectral properties as seen on the surfaces of the dark Saturnian moons—Iapetus and Hyperion. The most probable candidate is nanophase iron or hematite. Measurements of the thermal emission as well as radar observations by the Cassini spacecraft show that the icy regolith on the surface of Tethys is structurally complex and has a large porosity exceeding 95%. Surface features Color patterns The surface of Tethys has a number of large-scale features distinguished by their color and sometimes brightness. The trailing hemisphere gets increasingly red and dark as the anti-apex of motion is approached. This darkening is responsible for the hemispheric albedo asymmetry mentioned above. The leading hemisphere also reddens slightly as the apex of the motion is approached, although without any noticeable darkening. Such a bifurcated color pattern results in the existence of a bluish band between hemispheres following a great circle that runs through the poles. This coloration and darkening of the Tethyan surface is typical for Saturnian middle-sized satellites. Its origin may be related to a deposition of bright ice particles from the E-ring onto the leading hemispheres and dark particles coming from outer satellites on the trailing hemispheres. The darkening of the trailing hemispheres can also be caused by the impact of plasma from the magnetosphere of Saturn, which co-rotates with the planet. On the leading hemisphere of Tethys spacecraft observations have found a dark bluish band spanning 20° to the south and north from the equator. The band has an elliptical shape getting narrower as it approaches the trailing hemisphere. A comparable band exists only on Mimas. The band is almost certainly caused by the influence of energetic electrons from the Saturnian magnetosphere with energies greater than about 1 MeV. These particles drift in the direction opposite to the rotation of the planet and preferentially impact areas on the leading hemisphere close to the equator. Temperature maps of Tethys obtained by Cassini have shown this bluish region is cooler at midday than surrounding areas, giving the satellite a "Pac-man"-like appearance at mid-infrared wavelengths. Geology The surface of Tethys mostly consists of hilly cratered terrain dominated by craters more than 40 km in diameter. A smaller portion of the surface is represented by the smooth plains on the trailing hemisphere. There are also a number of tectonic features such as chasmata and troughs. The western part of the leading hemisphere of Tethys is dominated by Odysseus, a large impact basin whose 450 km diameter is nearly 2/5 of that of Tethys itself. The crater is now quite flat – more precisely, its floor conforms to Tethys's spherical shape. This is most likely due to the viscous relaxation of the Tethyan icy crust over geologic time. Nevertheless, the rim crest of Odysseus is elevated by approximately 5 km above the mean satellite radius. The central complex of Odysseus features a central pit 2–4 km deep surrounded by massifs elevated by 6–9 km above the crater floor, which itself is about 3 km below the average radius. The second major feature seen on Tethys is a huge valley called Ithaca Chasma, about 100 km wide and 3 km deep. It is more than 2,000 km in length, approximately 3/4 of the way around Tethys's circumference. Ithaca Chasma occupies about 10% of the surface of Tethys. It is approximately concentric with Odysseus—a pole of Ithaca Chasma lies only approximately 20° from the crater. It is thought that Ithaca Chasma formed as Tethys's internal liquid water solidified, causing the moon to expand and cracking the surface to accommodate the extra volume within. The subsurface ocean may have resulted from a 2:3 orbital resonance between Dione and Tethys early in the Solar System's history that led to orbital eccentricity and tidal heating of Tethys's interior. The ocean would have frozen after the moons escaped from the resonance. There is another theory about the formation of Ithaca Chasma: when the impact that caused the great crater Odysseus occurred, the shock wave traveled through Tethys and fractured the icy, brittle surface. In this case Ithaca Chasma would be the outermost ring graben of Odysseus. However, age determination based on crater counts in high-resolution Cassini images showed that Ithaca Chasma is older than Odysseus making the impact hypothesis unlikely. The smooth plains on the trailing hemisphere are approximately antipodal to Odysseus, although they extend about 60° to the northeast from the exact antipode. The plains have a relatively sharp boundary with the surrounding cratered terrain. The location of this unit near Odysseus' antipode argues for a connection between the crater and plains. The latter can be a result of focusing the seismic waves produced by the impact in the center of the opposite hemisphere. However the smooth appearance of the plains together with their sharp boundaries (impact shaking would have produced a wide transitional zone) indicates that they formed by endogenic intrusion, possibly along the lines of weakness in the Tethyan lithosphere created by Odysseus impact. Impact craters and chronology The majority of Tethyan impact craters are of a simple central peak type. Those more than 150 km in diameter show more complex peak ring morphology. Only Odysseus crater has a central depression resembling a central pit. Older impact craters are somewhat shallower than young ones implying a degree of relaxation. The density of impact craters varies across the surface of Tethys. The higher the crater density, the older the surface. This allows scientists to establish a relative chronology for Tethys. The cratered terrain is the oldest unit likely dating back to the Solar System formation 4.56 billion years ago. The youngest unit lies within Odysseus crater with an estimated age from 3.76 to 1.06 billion years, depending on the absolute chronology used. Ithaca Chasma is older than Odysseus. Origin and evolution Tethys is thought to have formed from an accretion disc or subnebula; a disc of gas and dust that existed around Saturn for some time after its formation. The low temperature at the position of Saturn in the Solar nebular means that water ice was the primary solid from which all moons formed. Other more volatile compounds like ammonia and carbon dioxide were likely present as well, though their abundances are not well constrained. The extremely water-ice-rich composition of Tethys remains unexplained. The conditions in the Saturnian sub-nebula likely favored conversion of the molecular nitrogen and carbon monoxide into ammonia and methane, respectively. This can partially explain why Saturnian moons including Tethys contain more water ice than outer Solar System bodies like Pluto or Triton as the oxygen freed from carbon monoxide would react with the hydrogen forming water. One of the most interesting explanations proposed is that the rings and inner moons accreted from the tidally stripped ice-rich crust of a Titan-like moon before it was swallowed by Saturn. The accretion process probably lasted for several thousand years before the moon was fully formed. Models suggest that impacts accompanying accretion caused heating of Tethys's outer layer, reaching a maximum temperature of around 155 K at a depth of about 29 km. After the end of formation due to thermal conduction, the subsurface layer cooled and the interior heated up. The cooling near-surface layer contracted and the interior expanded. This caused strong extensional stresses in Tethys's crust reaching estimates of 5.7 MPa, which likely led to cracking. Because Tethys lacks substantial rock content, the heating by decay of radioactive elements is unlikely to have played a significant role in its further evolution. This also means that Tethys may have never experienced any significant melting unless its interior was heated by tides. They may have occurred, for instance, during the passage of Tethys through an orbital resonance with Dione or some other moon. Still, present knowledge of the evolution of Tethys is very limited. Exploration Pioneer 11 flew by Saturn in 1979, and its closest approach to Tethys was 329,197 km on 1 September 1979. One year later, on 12 November 1980, Voyager 1 flew 415,670 km from Tethys. Its twin spacecraft, Voyager 2, passed as close as 93,010 km from the moon on 26 August 1981. Although both spacecraft took images of Tethys, the resolution of Voyager 1'''s images did not exceed 15 km, and only those obtained by Voyager 2 had a resolution as high as 2 km. The first geological feature discovered in 1980 by Voyager 1 was Ithaca Chasma. Later in 1981 Voyager 2 revealed that it almost circled the moon running for 270°. Voyager 2 also discovered the Odysseus crater. Tethys was the Saturnian satellite most fully imaged by the Voyagers. The Cassini spacecraft entered orbit around Saturn in 2004. During its primary mission from June 2004 through June 2008 it performed one very close targeted flyby of Tethys on 24 September 2005 at the distance of 1,503 km. In addition to this flyby the spacecraft performed many non-targeted flybys during its primary and equinox missions since 2004, at distances of tens of thousands of kilometers. Another flyby of Tethys took place on 14 August 2010 (during the solstice mission) at a distance of 38,300 km, when the fourth-largest crater on Tethys, Penelope, which is 207 km wide, was imaged. More non-targeted flybys were planned for the solstice mission in 2011–2017.Cassinis observations allowed high-resolution maps of Tethys to be produced with the resolution of 0.29 km. The spacecraft obtained spatially resolved near-infrared spectra of Tethys showing that its surface is made of water ice mixed with a dark material, whereas the far-infrared observations constrained the bolometric bond albedo. The radar observations at the wavelength of 2.2 cm showed that the ice regolith has a complex structure and is very porous. The observations of plasma in the vicinity of Tethys demonstrated that it is a geologically dead body producing no new plasma in the Saturnian magnetosphere. Future missions to Tethys and the Saturn system are uncertain, but one possibility is the Titan Saturn System Mission. Quadrangles Tethys is divided into 15 quadrangles: North Polar Area Anticleia Odysseus Alcinous Telemachus Circe Polycaste Theoclymenus Penelope Salmoneus Ithaca Chasma Hermione Melanthius Antinous South Polar Area Tethys in fiction	Physical sciences	Solar System	Astronomy
84726	https://en.wikipedia.org/wiki/Sulfate	Sulfate	The sulfate or sulphate ion is a polyatomic anion with the empirical formula . Salts, acid derivatives, and peroxides of sulfate are widely used in industry. Sulfates occur widely in everyday life. Sulfates are salts of sulfuric acid and many are prepared from that acid. Spelling "Sulfate" is the spelling recommended by IUPAC, but "sulphate" was traditionally used in British English. Structure The sulfate anion consists of a central sulfur atom surrounded by four equivalent oxygen atoms in a tetrahedral arrangement. The symmetry of the isolated anion is the same as that of methane. The sulfur atom is in the +6 oxidation state while the four oxygen atoms are each in the −2 state. The sulfate ion carries an overall charge of −2 and it is the conjugate base of the bisulfate (or hydrogensulfate) ion, , which is in turn the conjugate base of , sulfuric acid. Organic sulfate esters, such as dimethyl sulfate, are covalent compounds and esters of sulfuric acid. The tetrahedral molecular geometry of the sulfate ion is as predicted by VSEPR theory. Bonding The first description of the bonding in modern terms was by Gilbert Lewis in his groundbreaking paper of 1916 where he described the bonding in terms of electron octets around each atom, that is no double bonds and a formal charge of +2 on the sulfur atom and -1 on each oxygen atom. Later, Linus Pauling used valence bond theory to propose that the most significant resonance canonicals had two pi bonds involving d orbitals. His reasoning was that the charge on sulfur was thus reduced, in accordance with his principle of electroneutrality. The S−O bond length of 149 pm is shorter than the bond lengths in sulfuric acid of 157 pm for S−OH. The double bonding was taken by Pauling to account for the shortness of the S−O bond. Pauling's use of d orbitals provoked a debate on the relative importance of pi bonding and bond polarity (electrostatic attraction) in causing the shortening of the S−O bond. The outcome was a broad consensus that d orbitals play a role, but are not as significant as Pauling had believed. A widely accepted description involving pπ – dπ bonding was initially proposed by Durward William John Cruickshank. In this model, fully occupied p orbitals on oxygen overlap with empty sulfur d orbitals (principally the dz2 and dx2–y2). However, in this description, despite there being some π character to the S−O bonds, the bond has significant ionic character. For sulfuric acid, computational analysis (with natural bond orbitals) confirms a clear positive charge on sulfur (theoretically +2.45) and a low 3d occupancy. Therefore, the representation with four single bonds is the optimal Lewis structure rather than the one with two double bonds (thus the Lewis model, not the Pauling model). In this model, the structure obeys the octet rule and the charge distribution is in agreement with the electronegativity of the atoms. The discrepancy between the S−O bond length in the sulfate ion and the S−OH bond length in sulfuric acid is explained by donation of p-orbital electrons from the terminal S=O bonds in sulfuric acid into the antibonding S−OH orbitals, weakening them resulting in the longer bond length of the latter. However, Pauling's representation for sulfate and other main group compounds with oxygen is still a common way of representing the bonding in many textbooks. The apparent contradiction can be clarified if one realizes that the covalent double bonds in the Lewis structure actually represent bonds that are strongly polarized by more than 90% towards the oxygen atom. On the other hand, in the structure with a dipolar bond, the charge is localized as a lone pair on the oxygen. Preparation Typically metal sulfates are prepared by treating metal oxides, metal carbonates, or the metal itself with sulfuric acid: Although written with simple anhydrous formulas, these conversions generally are conducted in the presence of water. Consequently the product sulfates are hydrated, corresponding to zinc sulfate , copper(II) sulfate , and cadmium sulfate . Some metal sulfides can be oxidized to give metal sulfates. Properties There are numerous examples of ionic sulfates, many of which are highly soluble in water. Exceptions include calcium sulfate, strontium sulfate, lead(II) sulfate, barium sulfate, silver sulfate, and mercury sulfate, which are poorly soluble. Radium sulfate is the most insoluble sulfate known. The barium derivative is useful in the gravimetric analysis of sulfate: if one adds a solution of most barium salts, for instance barium chloride, to a solution containing sulfate ions, barium sulfate will precipitate out of solution as a whitish powder. This is a common laboratory test to determine if sulfate anions are present. The sulfate ion can act as a ligand attaching either by one oxygen (monodentate) or by two oxygens as either a chelate or a bridge. An example is the complex or the neutral metal complex where the sulfate ion is acting as a bidentate ligand. The metal–oxygen bonds in sulfate complexes can have significant covalent character. Uses and occurrence Commercial applications Sulfates are widely used industrially. Major compounds include: Gypsum, the natural mineral form of hydrated calcium sulfate, is used to produce plaster. About 100 million tonnes per year are used by the construction industry. Copper sulfate, a common algaecide, the more stable form () is used for galvanic cells as electrolyte Iron(II) sulfate, a common form of iron in mineral supplements for humans, animals, and soil for plants Magnesium sulfate (commonly known as Epsom salts), used in therapeutic baths Lead(II) sulfate, produced on both plates during the discharge of a lead–acid battery Sodium laureth sulfate, or SLES, a common detergent in shampoo formulations Polyhalite, , used as fertiliser. Occurrence in nature Sulfate-reducing bacteria, some anaerobic microorganisms, such as those living in sediment or near deep sea thermal vents, use the reduction of sulfates coupled with the oxidation of organic compounds or hydrogen as an energy source for chemosynthesis. History Some sulfates were known to alchemists. The vitriol salts, from the Latin vitreolum, glassy, were so-called because they were some of the first transparent crystals known. Green vitriol is iron(II) sulfate heptahydrate, ; blue vitriol is copper(II) sulfate pentahydrate, and white vitriol is zinc sulfate heptahydrate, . Alum, a double sulfate of potassium and aluminium with the formula , figured in the development of the chemical industry. Environmental effects Sulfates occur as microscopic particles (aerosols) resulting from fossil fuel and biomass combustion. They increase the acidity of the atmosphere and form acid rain. The anaerobic sulfate-reducing bacteria Desulfovibrio desulfuricans and D. vulgaris can remove the black sulfate crust that often tarnishes buildings. Main effects on climate Reversal and accelerated warming Hydrological cycle Solar geoengineering Hydrogensulfate (bisulfate) The hydrogensulfate ion (), also called the bisulfate ion, is the conjugate base of sulfuric acid (). Sulfuric acid is classified as a strong acid; in aqueous solutions it ionizes completely to form hydronium () and hydrogensulfate () ions. In other words, the sulfuric acid behaves as a Brønsted–Lowry acid and is deprotonated to form hydrogensulfate ion. Hydrogensulfate has a valency of 1. An example of a salt containing the ion is sodium bisulfate, . In dilute solutions the hydrogensulfate ions also dissociate, forming more hydronium ions and sulfate ions (). Other sulfur oxyanions	Physical sciences	Salts	null
84777	https://en.wikipedia.org/wiki/Fingerprint	Fingerprint	A fingerprint is an impression left by the friction ridges of a human finger. The recovery of partial fingerprints from a crime scene is an important method of forensic science. Moisture and grease on a finger result in fingerprints on surfaces such as glass or metal. Deliberate impressions of entire fingerprints can be obtained by ink or other substances transferred from the peaks of friction ridges on the skin to a smooth surface such as paper. Fingerprint records normally contain impressions from the pad on the last joint of fingers and thumbs, though fingerprint cards also typically record portions of lower joint areas of the fingers. Human fingerprints are detailed, unique, difficult to alter, and durable over the life of an individual, making them suitable as long-term markers of human identity. They may be employed by police or other authorities to identify individuals who wish to conceal their identity, or to identify people who are incapacitated or deceased and thus unable to identify themselves, as in the aftermath of a natural disaster. Their use as evidence has been challenged by academics, judges and the media. There are no uniform standards for point-counting methods, and academics have argued that the error rate in matching fingerprints has not been adequately studied and that fingerprint evidence has no secure statistical foundation. Research has been conducted into whether experts can objectively focus on feature information in fingerprints without being misled by extraneous information, such as context. Biology Fingerprints are impressions left on surfaces by the friction ridges on the finger of a human. The matching of two fingerprints is among the most widely used and most reliable biometric techniques. Fingerprint matching considers only the obvious features of a fingerprint. The composition of fingerprints consists of water (95%-99%), as well as organic and inorganic constituents. The organic component is made up of amino acids, proteins, glucose, lactase, urea, pyruvate, fatty acids and sterols. Inorganic ions such as chloride, sodium, potassium and iron are also present. Other contaminants such as oils found in cosmetics, drugs and their metabolites and food residues may be found in fingerprint residues. A friction ridge is a raised portion of the epidermis on the digits (fingers and toes), the palm of the hand or the sole of the foot, consisting of one or more connected ridge units of friction ridge skin. These are sometimes known as "epidermal ridges" which are caused by the underlying interface between the dermal papillae of the dermis and the interpapillary (rete) pegs of the epidermis. These unique features are formed at around the 15th week of fetal development and remain until after death, when decomposition begins. During the development of the fetus, around the 13th week of a pregnancy, ledge-like formation is formed at the bottom of the epidermis beside the dermis. The cells along these ledges begin to rapidly proliferate. This rapid proliferation forms primary and secondary ridges. Both the primary and secondary ridges act as a template for the outer layer of the skin to form the friction ridges seen on the surface of the skin. These epidermal ridges serve to amplify vibrations triggered, for example, when fingertips brush across an uneven surface, better transmitting the signals to sensory nerves involved in fine texture perception. These ridges may also assist in gripping rough surfaces and may improve surface contact in wet conditions. Genetics Consensus within the scientific community suggests that the dermatoglyphic patterns on fingertips are hereditary. The fingerprint patterns between monozygotic twins have been shown to be very similar (though not identical), whereas dizygotic twins have considerably less similarity. Significant heritability has been identified for 12 dermatoglyphic characteristics. Current models of dermatoglyphic trait inheritance suggest Mendelian transmission with additional effects from either additive or dominant major genes. Whereas genes determine the general characteristics of patterns and their type, the presence of environmental factors result in the slight differentiation of each fingerprint. However, the relative influences of genetic and environmental effects on fingerprint patterns are generally unclear. One study has suggested that roughly 5% of the total variability is due to small environmental effects, although this was only performed using total ridge count as a metric. Several models of finger ridge formation mechanisms that lead to the vast diversity of fingerprints have been proposed. One model suggests that a buckling instability in the basal cell layer of the fetal epidermis is responsible for developing epidermal ridges. Additionally, blood vessels and nerves may also serve a role in the formation of ridge configurations. Another model indicates that changes in amniotic fluid surrounding each developing finger within the uterus cause corresponding cells on each fingerprint to grow in different microenvironments. For a given individual, these various factors affect each finger differently, preventing two fingerprints from being identical while still retaining similar patterns. It is important to note that the determination of fingerprint inheritance is made difficult by the vast diversity of phenotypes. Classification of a specific pattern is often subjective (lack of consensus on the most appropriate characteristic to measure quantitatively) which complicates analysis of dermatoglyphic patterns. Several modes of inheritance have been suggested and observed for various fingerprint patterns. Total fingerprint ridge count, a commonly used metric of fingerprint pattern size, has been suggested to have a polygenic mode of inheritance and is influenced by multiple additive genes. This hypothesis has been challenged by other research, however, which indicates that ridge counts on individual fingers are genetically independent and lack evidence to support the existence of additive genes influencing pattern formation. Another mode of fingerprint pattern inheritance suggests that the arch pattern on the thumb and on other fingers are inherited as an autosomal dominant trait. Further research on the arch pattern has suggested that a major gene or multifactorial inheritance is responsible for arch pattern heritability. A separate model for the development of the whorl pattern indicates that a single gene or group of linked genes contributes to its inheritance. Furthermore, inheritance of the whorl pattern does not appear to be symmetric in that the pattern is seemingly randomly distributed among the ten fingers of a given individual. In general, comparison of fingerprint patterns between left and right hands suggests an asymmetry in the effects of genes on fingerprint patterns, although this observation requires further analysis. In addition to proposed models of inheritance, specific genes have been implicated as factors in fingertip pattern formation (their exact mechanism of influencing patterns is still under research). Multivariate linkage analysis of finger ridge counts on individual fingers revealed linkage to chromosome 5q14.1 specifically for the ring, index, and middle fingers. In mice, variants in the gene EVI1 were correlated with dermatoglyphic patterns. EVI1 expression in humans does not directly influence fingerprint patterns but does affect limb and digit formation which in turn may play a role in influencing fingerprint patterns. Genome-wide association studies found single nucleotide polymorphisms within the gene ADAMTS9-AS2 on 3p14.1, which appeared to have an influence on the whorl pattern on all digits. This gene encodes antisense RNA which may inhibit ADAMTS9, which is expressed in the skin. A model of how genetic variants of ADAMTS9-AS2 directly influence whorl development has not yet been proposed. In February 2023, a study identified WNT, BMP and EDAR as signaling pathways regulating the formation of primary ridges on fingerprints, with the first two having an opposite relationship established by a Turing reaction-diffusion system. Classification systems Before computerization, manual filing systems were used in large fingerprint repositories. A fingerprint classification system groups fingerprints according to their characteristics and therefore helps in the matching of a fingerprint against a large database of fingerprints. A query fingerprint that needs to be matched can therefore be compared with a subset of fingerprints in an existing database. Early classification systems were based on the general ridge patterns, including the presence or absence of circular patterns, of several or all fingers. This allowed the filing and retrieval of paper records in large collections based on friction ridge patterns alone. The most popular systems used the pattern class of each finger to form a numeric key to assist lookup in a filing system. Fingerprint classification systems included the Roscher System, the Juan Vucetich System and the Henry Classification System. The Roscher System was developed in Germany and implemented in both Germany and Japan. The Vucetich System was developed in Argentina and implemented throughout South America. The Henry Classification System was developed in India and implemented in most English-speaking countries. In the Henry Classification System, there are three basic fingerprint patterns: loop, whorl, and arch, which constitute 60–65 percent, 30–35 percent, and 5 percent of all fingerprints respectively. There are also more complex classification systems that break down patterns even further, into plain arches or tented arches, and into loops that may be radial or ulnar, depending on the side of the hand toward which the tail points. Ulnar loops start on the pinky-side of the finger, the side closer to the ulna, the lower arm bone. Radial loops start on the thumb-side of the finger, the side closer to the radius. Whorls may also have sub-group classifications including plain whorls, accidental whorls, double loop whorls, peacock's eye, composite, and central pocket loop whorls. The "primary classification number" in the Henry Classification System is a fraction whose numerator and denominator are whole numbers between 1 and 32 inclusive, thus classifying each set of ten fingerprints into one of 1024 groups. (To distinguish these groups, the fraction is not reduced by dividing out any common factors.) The fraction is determined by ten indicators, one for each finger, an indicator taking the value 1 when that finger has a whorl, and 0 otherwise. These indicators can be written for the right hand and for the left hand, where the subscripts are t for thumb, i for index finger, m for middle finger, r for ring finger and l for little finger. The formula for the fraction is then as follows: For example, if only the right ring finger and the left index finger have whorls, then the set of fingerprints is classified into the "9/3" group: Note that although 9/3 = 3/1, the "9/3" group is different from the "3/1" group, as the latter corresponds to having whorls only on the left middle finger. Fingerprint identification Fingerprint identification, known as dactyloscopy, ridgeology, or hand print identification, is the process of comparing two instances of friction ridge skin impressions (see minutiae), from human fingers or toes, or even the palm of the hand or sole of the foot, to determine whether these impressions could have come from the same individual. The flexibility and the randomized formation of the friction ridges on skin means that no two finger or palm prints are ever exactly alike in every detail; even two impressions recorded immediately after each other from the same hand may be slightly different. Fingerprint identification, also referred to as individualization, involves an expert, or an expert computer system operating under threshold scoring rules, determining whether two friction ridge impressions are likely to have originated from the same finger or palm (or toe or sole). In 2024, research using deep learning neural networks found contrary to "prevailing assumptions" that fingerprints from different fingers of the same person could be identified as belonging to that individual with 99.99% confidence. Further, features used in traditional methods were nonpredictive in such identification while ridge orientation, particularly near the center of the fingerprint center provided most information. An intentional recording of friction ridges is usually made with black printer's ink rolled across a contrasting white background, typically a white card. Friction ridges can also be recorded digitally, usually on a glass plate, using a technique called live scan. A "latent print" is the chance recording of friction ridges deposited on the surface of an object or a wall. Latent prints are invisible to the naked eye, whereas "patent prints" or "plastic prints" are viewable with the unaided eye. Latent prints are often fragmentary and require the use of chemical methods, powder, or alternative light sources in order to be made clear. Sometimes an ordinary bright flashlight will make a latent print visible. When friction ridges come into contact with a surface that will take a print, material that is on the friction ridges such as perspiration, oil, grease, ink, or blood, will be transferred to the surface. Factors which affect the quality of friction ridge impressions are numerous. Pliability of the skin, deposition pressure, slippage, the material from which the surface is made, the roughness of the surface, and the substance deposited are just some of the various factors which can cause a latent print to appear differently from any known recording of the same friction ridges. Indeed, the conditions surrounding every instance of friction ridge deposition are unique and never duplicated. For these reasons, fingerprint examiners are required to undergo extensive training. The scientific study of fingerprints is called dermatoglyphics or dactylography. Fingerprinting techniques Exemplar Exemplar prints, or known prints, is the name given to fingerprints deliberately collected from a subject, whether for purposes of enrollment in a system or when under arrest for a suspected criminal offense. During criminal arrests, a set of exemplar prints will normally include one print taken from each finger that has been rolled from one edge of the nail to the other, plain (or slap) impressions of each of the four fingers of each hand, and plain impressions of each thumb. Exemplar prints can be collected using live scan or by using ink on paper cards. Latent In forensic science, a partial fingerprint lifted from a surface is called a latent fingerprint. Moisture and grease on fingers result in latent fingerprints on surfaces such as glass. But because they are not clearly visible, their detection may require chemical development through powder dusting, the spraying of ninhydrin, iodine fuming, or soaking in silver nitrate. Depending on the surface or the material on which a latent fingerprint has been found, different methods of chemical development must be used. Forensic scientists use different techniques for porous surfaces, such as paper, and nonporous surfaces, such as glass, metal or plastic. Nonporous surfaces require the dusting process, where fine powder and a brush are used, followed by the application of transparent tape to lift the latent fingerprint off the surface. While the police often describe all partial fingerprints found at a crime scene as latent prints, forensic scientists call partial fingerprints that are readily visible patent prints. Chocolate, toner, paint or ink on fingers will result in patent fingerprints. Latent fingerprints impressions that are found on soft material, such as soap, cement or plaster, are called plastic prints by forensic scientists. Capture and detection Live scan devices Fingerprint image acquisition is considered to be the most critical step in an automated fingerprint authentication system, as it determines the final fingerprint image quality, which has a drastic effect on the overall system performance. There are different types of fingerprint readers on the market, but the basic idea behind each is to measure the physical difference between ridges and valleys. All the proposed methods can be grouped into two major families: solid-state fingerprint readers and optical fingerprint readers. The procedure for capturing a fingerprint using a sensor consists of rolling or touching with the finger onto a sensing area, which according to the physical principle in use (optical, ultrasonic, capacitive, or thermalsee ) captures the difference between valleys and ridges. When a finger touches or rolls onto a surface, the elastic skin deforms. The quantity and direction of the pressure applied by the user, the skin conditions and the projection of an irregular 3D object (the finger) onto a 2D flat plane introduce distortions, noise, and inconsistencies in the captured fingerprint image. These problems result in inconsistent and non-uniform irregularities in the image. During each acquisition, therefore, the results of the imaging are different and uncontrollable. The representation of the same fingerprint changes every time the finger is placed on the sensor plate, increasing the complexity of any attempt to match fingerprints, impairing the system performance and consequently, limiting the widespread use of this biometric technology. In order to overcome these problems, as of 2010, non-contact or touchless 3D fingerprint scanners have been developed. Acquiring detailed 3D information, 3D fingerprint scanners take a digital approach to the analog process of pressing or rolling the finger. By modelling the distance between neighboring points, the fingerprint can be imaged at a resolution high enough to record all the necessary detail. Fingerprinting on cadavers The human skin itself, which is a regenerating organ until death, and environmental factors such as lotions and cosmetics, pose challenges when fingerprinting a human. Following the death of a human, the skin dries and cools. Fingerprints of dead humans may be obtained during an autopsy. The collection of fingerprints off of a cadaver can be done in varying ways and depends on the condition of the skin. In the case of cadaver in the later stages of decomposition with dried skin, analysts will boil the skin to recondition/rehydrate it, allowing for moisture to flow back into the skin and resulting in detail friction ridges. Another method that has been used in brushing a powder, such as baby powder over the tips of the fingers. The powder will ebbed itself into the farrows of the friction ridges allowing for the lifted ridges to be seen. Latent fingerprint detection In the 1930s, criminal investigators in the United States first discovered the existence of latent fingerprints on the surfaces of fabrics, most notably on the insides of gloves discarded by perpetrators. Since the late nineteenth century, fingerprint identification methods have been used by police agencies around the world to identify suspected criminals as well as the victims of crime. The basis of the traditional fingerprinting technique is simple. The skin on the palmar surface of the hands and feet forms ridges, so-called papillary ridges, in patterns that are unique to each individual and which do not change over time. Even identical twins (who share their DNA) do not have identical fingerprints. The best way to render latent fingerprints visible, so that they can be photographed, can be complex and may depend, for example, on the type of surfaces on which they have been left. It is generally necessary to use a "developer", usually a powder or chemical reagent, to produce a high degree of visual contrast between the ridge patterns and the surface on which a fingerprint has been deposited. Developing agents depend on the presence of organic materials or inorganic salts for their effectiveness, although the water deposited may also take a key role. Fingerprints are typically formed from the aqueous-based secretions of the eccrine glands of the fingers and palms with additional material from sebaceous glands primarily from the forehead. This latter contamination results from the common human behaviors of touching the face and hair. The resulting latent fingerprints consist usually of a substantial proportion of water with small traces of amino acids and chlorides mixed with a fatty, sebaceous component which contains a number of fatty acids and triglycerides. Detection of a small proportion of reactive organic substances such as urea and amino acids is far from easy. Fingerprints at a crime scene may be detected by simple powders, or by chemicals applied in situ. More complex techniques, usually involving chemicals, can be applied in specialist laboratories to appropriate articles removed from a crime scene. With advances in these more sophisticated techniques, some of the more advanced crime scene investigation services from around the world were, as of 2010, reporting that 50% or more of the fingerprints recovered from a crime scene had been identified as a result of laboratory-based techniques. Forensic laboratories Although there are hundreds of reported techniques for fingerprint detection, many of these are only of academic interest and there are only around 20 really effective methods which are currently in use in the more advanced fingerprint laboratories around the world. Some of these techniques, such as ninhydrin, diazafluorenone and vacuum metal deposition, show great sensitivity and are used operationally. Some fingerprint reagents are specific, for example ninhydrin or diazafluorenone reacting with amino acids. Others such as ethyl cyanoacrylate polymerisation, work apparently by water-based catalysis and polymer growth. Vacuum metal deposition using gold and zinc has been shown to be non-specific, but can detect fat layers as thin as one molecule. More mundane methods, such as the application of fine powders, work by adhesion to sebaceous deposits and possibly aqueous deposits in the case of fresh fingerprints. The aqueous component of a fingerprint, while initially sometimes making up over 90% of the weight of the fingerprint, can evaporate quite quickly and may have mostly gone after 24 hours. Following work on the use of argon ion lasers for fingerprint detection, a wide range of fluorescence techniques have been introduced, primarily for the enhancement of chemically developed fingerprints; the inherent fluorescence of some latent fingerprints may also be detected. Fingerprints can for example be visualized in 3D and without chemicals by the use of infrared lasers. A comprehensive manual of the operational methods of fingerprint enhancement was last published by the UK Home Office Scientific Development Branch in 2013 and is used widely around the world. A technique proposed in 2007 aims to identify an individual's ethnicity, sex, and dietary patterns. Limitations and implications in a forensic context One of the main limitations of friction ridge impression evidence regarding the actual collection would be the surface environment, specifically talking about how porous the surface the impression is on. With non-porous surfaces, the residues of the impression will not be absorbed into the material of the surface, but could be smudged by another surface. With porous surfaces, the residues of the impression will be absorbed into the surface. With both resulting in either an impression of no value to examiners or the destruction of the friction ridge impressions. In order for analysts to correctly positively identify friction ridge patterns and their features depends heavily on the clarity of the impression. Therefore, the analysis of friction ridges is limited by clarity. In a court context, many have argued that friction ridge identification and ridgeology should be classified as opinion evidence and not as fact, therefore should be assessed as such. Many have said that friction ridge identification is only legally admissible today because during the time when it was added to the legal system, the admissibility standards were quite low. There are only a limited number of studies that have been conducted to help confirm the science behind this identification process. Crime scene investigations The application of the new scanning Kelvin probe (SKP) fingerprinting technique, which makes no physical contact with the fingerprint and does not require the use of developers, has the potential to allow fingerprints to be recorded while still leaving intact material that could subsequently be subjected to DNA analysis. A forensically usable prototype was under development at Swansea University during 2010, in research that was generating significant interest from the British Home Office and a number of different police forces across the UK, as well as internationally. The hope is that this instrument could eventually be manufactured in sufficiently large numbers to be widely used by forensic teams worldwide. Detection of drug use The secretions, skin oils and dead cells in a human fingerprint contain residues of various chemicals and their metabolites present in the body. These can be detected and used for forensic purposes. For example, the fingerprints of tobacco smokers contain traces of cotinine, a nicotine metabolite; they also contain traces of nicotine itself. Caution should be used, as its presence may be caused by mere contact of the finger with a tobacco product. By treating the fingerprint with gold nanoparticles with attached cotinine antibodies, and then subsequently with a fluorescent agent attached to cotinine antibodies, the fingerprint of a smoker becomes fluorescent; non-smokers' fingerprints stay dark. The same approach, as of 2010, is being tested for use in identifying heavy coffee drinkers, cannabis smokers, and users of various other drugs. Police force databases Most American law enforcement agencies use Wavelet Scalar Quantization (WSQ), a wavelet-based system for efficient storage of compressed fingerprint images at 500 pixels per inch (ppi). WSQ was developed by the FBI, the Los Alamos National Lab, and the National Institute of Standards and Technology (NIST). For fingerprints recorded at 1000 ppi spatial resolution, law enforcement (including the FBI) uses JPEG 2000 instead of WSQ. Validity Fingerprints collected at a crime scene, or on items of evidence from a crime, have been used in forensic science to identify suspects, victims and other persons who touched a surface. Fingerprint identification emerged as an important system within police agencies in the late 19th century, when it replaced anthropometric measurements as a more reliable method for identifying persons having a prior record, often under a false name, in a criminal record repository. Fingerprinting has served all governments worldwide during the past 100 years or so to provide identification of criminals. Fingerprints are the fundamental tool in every police agency for the identification of people with a criminal history. The validity of forensic fingerprint evidence has been challenged by academics, judges and the media. In the United States fingerprint examiners have not developed uniform standards for the identification of an individual based on matching fingerprints. In some countries where fingerprints are also used in criminal investigations, fingerprint examiners are required to match a number of identification points before a match is accepted. In England 16 identification points are required and in France 12, to match two fingerprints and identify an individual. Point-counting methods have been challenged by some fingerprint examiners because they focus solely on the location of particular characteristics in fingerprints that are to be matched. Fingerprint examiners may also uphold the one dissimilarity doctrine, which holds that if there is one dissimilarity between two fingerprints, the fingerprints are not from the same finger. Furthermore, academics have argued that the error rate in matching fingerprints has not been adequately studied and it has even been argued that fingerprint evidence has no secure statistical foundation. Research has been conducted into whether experts can objectively focus on feature information in fingerprints without being misled by extraneous information, such as context. Fingerprints can theoretically be forged and planted at crime scenes. Professional certification Fingerprinting was the basis upon which the first forensic professional organization was formed, the International Association for Identification (IAI), in 1915. The first professional certification program for forensic scientists was established in 1977, the IAI's Certified Latent Print Examiner program, which issued certificates to those meeting stringent criteria and had the power to revoke certification where an individual's performance warranted it. Other forensic disciplines have followed suit and established their own certification programs. History Antiquity and the medieval period Fingerprints have been found on ancient clay tablets, seals, and pottery. They have also been found on the walls of Egyptian tombs and on Minoan, Greek, and Chinese pottery. In ancient China officials authenticated government documents with their fingerprints. In about 200 BC, fingerprints were used to sign written contracts in Babylon. Fingerprints from 3D-scans of cuneiform tablets are extracted using the GigaMesh Software Framework. With the advent of silk and paper in China, parties to a legal contract impressed their handprints on the document. Sometime before 851 CE, an Arab merchant in China, Abu Zayd Hasan, witnessed Chinese merchants using fingerprints to authenticate loans.	Biology and health sciences	Integumentary system	null
84936	https://en.wikipedia.org/wiki/Blood%E2%80%93brain%20barrier	Blood–brain barrier	The blood–brain barrier (BBB) is a highly selective semipermeable border of endothelial cells that regulates the transfer of solutes and chemicals between the circulatory system and the central nervous system, thus protecting the brain from harmful or unwanted substances in the blood. The blood–brain barrier is formed by endothelial cells of the capillary wall, astrocyte end-feet ensheathing the capillary, and pericytes embedded in the capillary basement membrane. This system allows the passage of some small molecules by passive diffusion, as well as the selective and active transport of various nutrients, ions, organic anions, and macromolecules such as glucose and amino acids that are crucial to neural function. The blood–brain barrier restricts the passage of pathogens, the diffusion of solutes in the blood, and large or hydrophilic molecules into the cerebrospinal fluid, while allowing the diffusion of hydrophobic molecules (O2, CO2, hormones) and small non-polar molecules. Cells of the barrier actively transport metabolic products such as glucose across the barrier using specific transport proteins. The barrier also restricts the passage of peripheral immune factors, like signaling molecules, antibodies, and immune cells, into the central nervous system, thus insulating the brain from damage due to peripheral immune events. Specialized brain structures participating in sensory and secretory integration within brain neural circuits—the circumventricular organs and choroid plexus—have in contrast highly permeable capillaries. Structure The BBB results from the selectivity of the tight junctions between the endothelial cells of brain capillaries, restricting the passage of solutes. At the interface between blood and the brain, endothelial cells are adjoined continuously by these tight junctions, which are composed of smaller subunits of transmembrane proteins, such as occludin, claudins (such as Claudin-5), junctional adhesion molecule (such as JAM-A). Each of these tight junction proteins is stabilized to the endothelial cell membrane by another protein complex that includes scaffolding proteins such as tight junction protein 1 (ZO1) and associated proteins. The BBB is composed of endothelial cells restricting passage of substances from the blood more selectively than endothelial cells of capillaries elsewhere in the body. Astrocyte cell projections called astrocytic feet (also known as "glia limitans") surround the endothelial cells of the BBB, providing biochemical support to those cells. The BBB is distinct from the quite similar blood-cerebrospinal fluid barrier, which is a function of the choroidal cells of the choroid plexus, and from the blood-retinal barrier, which can be considered a part of the whole realm of such barriers. Not all vessels in the human brain exhibit BBB properties. Some examples of this include the circumventricular organs, the roof of the third and fourth ventricles, capillaries in the pineal gland on the roof of the diencephalon and the pineal gland. The pineal gland secretes the hormone melatonin "directly into the systemic circulation", thus melatonin is not affected by the blood–brain barrier. Development The BBB appears to be functional by the time of birth. P-glycoprotein, a transporter, exists already in the embryonal endothelium. Measurement of brain uptake of various blood-borne solutes showed that newborn endothelial cells were functionally similar to those in adults, indicating that a selective BBB is operative at birth. In mice, Claudin-5 loss during development is lethal and results in size-selective (up to 742Da) loosening of the BBB. Mosaic deletion of claudin-5 in adult endothelial cells (in mice) reveals BBB leakage up to 10kDa molecule 6 days after deletion of claudin-5 and lethality after 10 days after deletion demonstrating a critical role of Claudin-5 in adult BBB. Function The blood–brain barrier acts effectively to protect brain tissue from circulating pathogens and other potentially toxic substances. Accordingly, blood-borne infections of the brain are rare. Infections of the brain that do occur are often difficult to treat. Antibodies are too large to cross the blood–brain barrier, and only certain antibiotics are able to pass. In some cases, a drug has to be administered directly into the cerebrospinal fluid where it can enter the brain by crossing the blood-cerebrospinal fluid barrier. Circumventricular organs Circumventricular organs (CVOs) are individual structures located adjacent to the fourth ventricle or third ventricle in the brain, and are characterized by dense capillary beds with permeable endothelial cells unlike those of the blood–brain barrier. Included among CVOs having highly permeable capillaries are the area postrema, subfornical organ, vascular organ of the lamina terminalis, median eminence, pineal gland, and three lobes of the pituitary gland. Permeable capillaries of the sensory CVOs (area postrema, subfornical organ, vascular organ of the lamina terminalis) enable rapid detection of circulating signals in systemic blood, while those of the secretory CVOs (median eminence, pineal gland, pituitary lobes) facilitate transport of brain-derived signals into the circulating blood. Consequently, the CVO permeable capillaries are the point of bidirectional blood–brain communication for neuroendocrine function. Specialized permeable zones The border zones between brain tissue "behind" the blood–brain barrier and zones "open" to blood signals in certain CVOs contain specialized hybrid capillaries that are leakier than typical brain capillaries, but not as permeable as CVO capillaries. Such zones exist at the border of the area postrema—nucleus tractus solitarii (NTS), and median eminence—hypothalamic arcuate nucleus. These zones appear to function as rapid transit regions for brain structures involved in diverse neural circuits—like the NTS and arcuate nucleus—to receive blood signals which are then transmitted into neural output. The permeable capillary zone shared between the median eminence and hypothalamic arcuate nucleus is augmented by wide pericapillary spaces, facilitating bidirectional flow of solutes between the two structures, and indicating that the median eminence is not only a secretory organ, but may also be a sensory organ. Therapeutic research As a drug target The blood–brain barrier is formed by the brain capillary endothelium and excludes from the brain 100% of large-molecule neurotherapeutics and more than 98% of all small-molecule drugs. Overcoming the difficulty of delivering therapeutic agents to specific regions of the brain presents a major challenge to treatment of most brain disorders. In its neuroprotective role, the blood–brain barrier functions to hinder the delivery of many potentially important diagnostic and therapeutic agents to the brain. Therapeutic molecules and antibodies that might otherwise be effective in diagnosis and therapy do not cross the BBB in adequate amounts to be clinically effective. To overcome this problem some peptides able to naturally cross the BBB have been widely investigated as a drug delivery system. Mechanisms for drug targeting in the brain involve going either "through" or "behind" the BBB. Modalities for drug delivery to the brain in unit doses through the BBB entail its disruption by osmotic means, or biochemically by the use of vasoactive substances, such as bradykinin, or even by localized exposure to high-intensity focused ultrasound (HIFU). Other methods used to get through the BBB may entail the use of endogenous transport systems, including carrier-mediated transporters, such as glucose and amino acid carriers, receptor-mediated transcytosis for insulin or transferrin, and the blocking of active efflux transporters such as p-glycoprotein. Some studies have shown that vectors targeting BBB transporters, such as the transferrin receptor, have been found to remain entrapped in brain endothelial cells of capillaries, instead of being ferried across the BBB into the targeted area. Intranasal administration The brain can be targeted non-invasively via the nasal passage. The drugs that remain in the passage after mucociliary clearance, enter the brain via three pathways: (1) Olfactory nerve-olfactory bulb-brain; (2) Trigeminal nerve-brain; and (3) Lungs/ Gastrointestinal tract-blood–brain The first and second methods involve the nerves, so they use the neuronal pathway and the third is via systemic circulation. However, these methods are less efficient to deliver drugs as they are indirect methods. Nanoparticles Nanotechnology is under preliminary research for its potential to facilitate the transfer of drugs across the BBB. Capillary endothelial cells and associated pericytes may be abnormal in tumors and the blood–brain barrier may not always be intact in brain tumors. Other factors, such as astrocytes, may contribute to the resistance of brain tumors to therapy using nanoparticles. Fat soluble molecules less than 400 daltons in mass can freely diffuse past the BBB through lipid mediated passive diffusion. Damage in injury and disease The blood–brain barrier may become damaged in certain neurological diseases, as indicated by neuroimaging studies of Alzheimer's disease, amyotrophic lateral sclerosis, epilepsy, ischemic stroke, and brain trauma, and in systemic diseases, such as liver failure. Effects such as impaired glucose transport and endothelial degeneration may lead to metabolic dysfunction within the brain, and an increased permeability of the BBB to proinflammatory factors, potentially allowing antibiotics and phagocytes to move across the BBB. However, in many neurodegenerative diseases, the exact cause and pathology remains unknown. It is still unclear whether the BBB dysfunction in the disease is a causative agent, a result of the disease, or somewhere in the middle. History A 1898 study observed that low-concentration "bile salts" failed to affect behavior when injected into the blood of animals. Thus, in theory, the salts failed to enter the brain. Two years later, Max Lewandowsky may have been the first to coin the term "blood–brain barrier" in 1900, referring to the hypothesized semipermeable membrane. There is some debate over the creation of the term blood–brain barrier as it is often attributed to Lewandowsky, but it does not appear in his papers. The creator of the term may have been Lina Stern. Stern was a Russian scientist who published her work in Russian and French. Due to the language barrier between her publications and English-speaking scientists, this could have made her work a lesser-known origin of the term. All the while, bacteriologist Paul Ehrlich was studying staining, a procedure that is used in many microscopy studies to make fine biological structures visible using chemical dyes. As Ehrlich injected some of these dyes (notably the aniline dyes that were then widely used), the dye stained all of the organs of some kinds of animals except for their brains. At that time, Ehrlich attributed this lack of staining to the brain simply not picking up as much of the dye. However, in a later experiment in 1913, Edwin Goldmann (one of Ehrlich's students) injected the dye directly into the cerebrospinal fluid of animal brains. He found then the brains did become dyed, but the rest of the body did not, demonstrating the existence of a compartmentalization between the two. At that time, it was thought that the blood vessels themselves were responsible for the barrier, since no obvious membrane could be found.	Biology and health sciences	Circulatory system	Biology
85012	https://en.wikipedia.org/wiki/Sewing%20machine	Sewing machine	A sewing machine is a machine used to sew fabric and materials together with thread. Sewing machines were invented during the first Industrial Revolution to decrease the amount of manual sewing work performed in clothing companies. Since the invention of the first sewing machine, generally considered to have been the work of Englishman Thomas Saint in 1790, the sewing machine has greatly improved the efficiency and productivity of the clothing industry. Home sewing machines are designed for one person to sew individual items while using a single stitch type at a time. In a modern sewing machine, the process of stitching has been automated, so that the fabric easily glides in and out of the machine. Early sewing machines were powered by either constantly turning a flywheel handle or with a foot-operated treadle mechanism. Electrically-powered machines were later introduced. Industrial sewing machines, by contrast to domestic machines, are larger, faster, and more varied in their size, cost, appearance, and tasks. History Invention Charles Fredrick Wiesenthal, a German-born engineer working in England, was awarded the first British patent for a mechanical device to aid the art of sewing, in 1755. His invention consisted of a double pointed needle with an eye at one end. In 1790, the English inventor Thomas Saint invented the first sewing machine design. His machine was meant to be used on leather and canvas material. It is likely that Saint had a working model, but there is no surviving evidence of one. He was a skilled cabinet maker and his device included many practical and functional features: an overhanging arm; a feed mechanism (adequate for short lengths of leather); a vertical needle bar; and a looper. Saint created the machine to reduce the amount of hand-stitching on garments, making sewing more reliable and functional. His sewing machine used the chain stitch method, in which the machine uses a single thread to make simple stitches in the fabric. A stitching awl would have pierced the material, and a forked-point rod would have carried the thread through the hole, where it would have been hooked underneath and moved to the next stitching place, after which the cycle would be repeated, thereby locking the stitch in place. Saint's machine was designed to aid in the manufacturing of various leather goods, including saddles and bridles, but it was also capable of working with canvas, and was used for sewing ship sails. Although his machine was very advanced for the era, the concept would need steady improvement over the coming decades before it was practical enough to enter into wide use. In 1874, a sewing machine manufacturer, William Newton Wilson, found Saint's drawings in the UK Patent Office, made adjustments to the looper, and built a working machine, currently owned by the Science Museum in London. In 1804, a sewing machine was built by two Englishmen, Thomas Stone and James Henderson, and a machine for embroidering was constructed by John Duncan in Scotland. An Austrian tailor, Josef Madersperger, began developing his first sewing machine in 1807 and presented his first working machine publicly in 1814. Having received financial support from his government, the Austrian tailor worked on the development of his machine until 1839, when he built a machine imitating the weaving process using the chain stitch. The first practical and widely used sewing machine was invented by Barthélemy Thimonnier, a French tailor, in 1829. His machine sewed straight seams using a chain stitch like Saint's model had, and in 1830, he signed a contract with Auguste Ferrand, a mining engineer, who made the requisite drawings and submitted a patent application. The patent for his machine was issued on 17 July 1830, and in the same year, he and his partners opened the first machine-based clothing manufacturing company in the world to create army uniforms for the French Army. However, the factory was burned down, reportedly by workers fearful of losing their livelihood, following the issuing of the patent. A model of the machine is exhibited in London at the Science Museum. The machine is made of wood and uses a barbed needle which passes downward through the cloth to grab the thread and pull it up to form a loop to be locked by the next loop. The first American lockstitch sewing machine was invented by Walter Hunt in 1832. His machine used a needle with the eye and the point on the same end carrying the upper thread, and a falling shuttle carrying the lower thread. The curved needle moved through the fabric horizontally, leaving the loop as it withdrew. The shuttle passed through the loop, interlocking the thread. The feed was unreliable, requiring the machine to be stopped frequently and reset up. Hunt eventually lost interest in his machine and sold individual machines without bothering to patent his invention, and only patenting it at a late date of 1854. In 1842, John Greenough patented the first sewing machine in the United States. The British partners Newton and Archibold introduced the eye-pointed needle and the use of two pressing surfaces to keep the pieces of fabric in position, in 1841. The first machine to combine all the disparate elements of the previous half-century of innovation into the modern sewing machine was the device built by English inventor John Fisher in 1844, a little earlier than the very similar machines built by Isaac Merritt Singer in 1851, and the lesser known Elias Howe, in 1845. However, due to the botched filing of Fisher's patent at the Patent Office, he did not receive due recognition for the modern sewing machine in the legal disputations of priority with Singer, and Singer reaped the benefits of the patent. Industrial competition Elias Howe, born in Spencer, Massachusetts, created his sewing machine in 1845, using a similar method to Fisher's except that the fabric was held vertically. An important improvement on his machine was to have the needle running away from the point, starting from the eye. After a lengthy stay in England trying to attract interest in his machine, he returned to America to find various people infringing his patent, among them Isaac Merritt Singer. He eventually won a case for patent infringement in 1854 and was awarded the right to claim royalties from the manufacturers using ideas covered by his patent, including Singer. Singer had seen a rotary sewing machine being repaired in a Boston shop. As an engineer, he thought it was clumsy and decided to design a better one. The machine he devised used a falling shuttle instead of a rotary one; the needle was mounted vertically and included a presser foot to hold the cloth in place. It had a fixed arm to hold the needle and included a basic tension system. This machine combined elements of Thimonnier, Hunt and Howe's machines. Singer was granted an American patent in 1851. The foot treadle used since the Middle Ages, used to convert reciprocating to rotary motion, was adapted to drive the sewing machine, leaving both hands free. When Howe learned of Singer's machine he took him to court, where Howe won and Singer was forced to pay a lump sum for all machines already produced. Singer then took out a license under Howe's patent and paid him US$1.15 per machine before entering into a joint partnership with a lawyer named Edward Clark. They created the first hire-purchase arrangement to allow people to purchase their machines through payments over time. Meanwhile, Allen B. Wilson developed a shuttle that reciprocated in a short arc, which was an improvement over Singer and Howe's. However, John Bradshaw had patented a similar device and threatened to sue, so Wilson decided to try a new method. He went into partnership with Nathaniel Wheeler to produce a machine with a rotary hook instead of a shuttle. This was far quieter and smoother than other methods, with the result that the Wheeler & Wilson Company produced more machines in the 1850s and 1860s than any other manufacturer. Wilson also invented the four-motion feed mechanism that is still used on every sewing machine today. This had a forward, down, back and up motion, which drew the cloth through in an even and smooth motion. Charles Miller patented the first machine to stitch buttonholes. Throughout the 1850s more and more companies were being formed, each trying to sue the others for patent infringement. This triggered a patent thicket known as the Sewing Machine War. In 1856, the Sewing Machine Combination was formed, consisting of Singer, Howe, Wheeler, Wilson, and Grover and Baker. These four companies pooled their patents, with the result that all other manufacturers had to obtain a license for $15 per machine. This lasted until 1877 when the last patent expired. James Edward Allen Gibbs (1829–1902), a farmer from Raphine in Rockbridge County, Virginia, patented the first chain stitch single-thread sewing machine on June 2, 1857. In partnership with James Willcox, Gibbs became a principal partner in Willcox & Gibbs Sewing Machine Company. Willcox & Gibbs commercial sewing machines are still used in the 21st century, with spares parts available. Market expansion William Jones started making sewing machines in 1859 and in 1860 formed a partnership with Thomas Chadwick. As Chadwick & Jones, they manufactured sewing machines at Ashton-under-Lyne, England until 1863. Their machines used designs from Howe and Wilson produced under licence. Thomas Chadwick later joined Bradbury & Co. William Jones opened a factory in Guide Bridge, Manchester in 1869. In 1893 a Jones advertising sheet claimed that this factory was the "Largest Factory in England Exclusively Making First Class Sewing Machines". The firm was renamed as the Jones Sewing Machine Co. Ltd and was later acquired by Brother Industries of Japan, in 1968. Clothing manufacturers were the first sewing machine customers, and used them to produce the first ready-to-wear clothing and shoes. In the 1860s consumers began purchasing them, and the machines—ranging in price from £6 to £15 in Britain depending on features—became very common in middle-class homes. Owners were much more likely to spend free time with their machines to make and mend clothing for their families than to visit friends, and women's magazines and household guides such as Mrs Beeton's offered dress patterns and instructions. A sewing machine could produce a man's shirt in about one hour, compared to hours by hand. In 1877, the world's first crochet machine was invented and patented by Joseph M. Merrow, then-president of what had started in the 1840s, as a machine shop to develop specialized machinery for the knitting operations. This crochet machine was the first production overlock sewing machine. The Merrow Machine Company went on to become one of the largest American manufacturers of overlock sewing machines and remains in the 21st century as the last American over-lock sewing machine manufacturer. In 1885 Singer patented the Singer Vibrating Shuttle sewing machine, which used Allen B. Wilson's idea for a vibrating shuttle and was a better lockstitcher than the oscillating shuttles of the time. Millions of the machines, perhaps the world's first really practical sewing machine for domestic use, were produced until finally superseded by rotary shuttle machines in the 20th century. Sewing machines continued being made to roughly the same design—with more lavish decoration—until well into the 1900s. The first electric machines were developed by Singer Sewing Co. and introduced in 1889. By the end of the First World War, Singer was offering hand, treadle and electric machines for sale. At first, the electric machines were standard machines with a motor strapped on the side, but as more homes gained power, they became more popular, and the motor was gradually introduced into the casing. Introduction of electronic machines Sewing machines were strictly mechanical, using gears, shafts, levers, and so on, until the 1970s when electronic machines were introduced to the market. Electronic sewing machines incorporate components such as circuit boards, computer chips, and additional motors for independent control of machine functions. These electronic components enabled new features such as automating thread cutters, needle positioning, and back-tacking, as well as digitized stitch patterns and stitch combinations. Because of the lifespan and increased complexity of the electronic parts, electronic sewing machines do not last as long as mechanical sewing machines, which can last over 100 years. Stitches Sewing machines can make a great variety of plain or patterned stitches. Ignoring strictly decorative aspects, over three dozen distinct stitch formations are formally recognized by the ISO 4915:1991 standard, involving one to seven separate threads to form the stitch. Plain stitches fall into four general categories: chainstitch, lockstitch, overlock, and coverstitch. Chain stitch was used by early sewing machines and has two major drawbacks: The stitch is not self-locking, and if the thread breaks at any point or is not tied at both ends, the whole length of stitching comes out. It is also easily ripped out. The direction of sewing cannot be changed much from one stitch to the next, or the stitching process fails. A better stitch was found in the lockstitch. The chain stitch is still used today in clothing manufacture, though due to its major drawbacks, it is generally paired with an overlock stitch along the same seam. Lockstitch is the familiar stitch performed by most household sewing machines and most industrial "single needle" sewing machines, using two threads, one passed through a needle and one coming from a bobbin or shuttle. Each thread stays on its own side of the material while being sewn, interlacing with the other thread at each needle hole by means of a bobbin driver. As a result, a lockstitch can be formed anywhere on the material being sewn; it does not need to be near an edge. Overlock, also known as "serging" or "serger stitch", can be formed with two to four threads, one or two needles, and one or two loopers. Overlock sewing machines are usually equipped with knives that trim or create the edge immediately in front of the stitch formation. Household and industrial overlock machines are commonly used for garment seams in knit or stretchy fabrics, for garment seams where the fabric is light enough that the seam does not need to be pressed open, and for protecting edges against raveling. Machines using two to four threads are most common, and frequently one machine can be configured for several varieties of overlock stitch. Overlock machines with five or more threads usually make both a chainstitch with one needle and one looper, and an overlock stitch with the remaining needles and loopers. This combination is known as a "safety stitch". A similar machine used for stretch fabrics is called a mock safety. Coverstitch is formed by two or more needles and one or two loopers. Like lockstitch and chainstitch, coverstitch can be formed anywhere on the material being sewn. One looper manipulates a thread below the material being sewn, forming a bottom cover stitch against the needle threads. An additional looper above the material can form a top cover stitch simultaneously. The needle threads form parallel rows, while the looper threads cross back and forth all the needle rows. Coverstitch is so-called because the grid of crossing needle and looper threads covers raw seam edges, much as the overlock stitch does. It is widely used in garment construction, particularly for attaching trims and flat seaming where the raw edges can be finished in the same operation as forming the seam. A zigzag stitch is a variant geometry of the lockstitch. It is a back-and-forth stitch used where a straight stitch will not suffice, such as in preventing raveling of a fabric, in stitching stretchable fabrics, and in temporarily joining two work pieces edge-to-edge. When creating a zigzag stitch, the back-and-forth motion of the sewing machine's needle is controlled by a cam. As the cam rotates, a fingerlike follower that is connected to the needle bar rides along the cam and tracks its indentations. As the follower moves in and out, the needle bar is moved from side to side. Very old sewing machines lack this hardware and so cannot natively produce a zigzag stitch, but there are often shank-driven attachments available which enable them to do so. Feed mechanisms Besides the basic motion of needles, loopers and bobbins, the material being sewn must move so that each cycle of needle motion involves a different part of the material. This motion is known as feed, and sewing machines have almost as many ways of feeding material as they do of forming stitches. For general categories, there are: drop feed, needle feed, walking foot, puller, and manual. Often, multiple types of feed are used on the same machine. Besides these general categories, there are also uncommon feed mechanisms used in specific applications like edge joining fur, making seams on caps, and blind stitching. The drop feed mechanism is used by almost all household machines and involves a mechanism below the sewing surface of the machine. When the needle is withdrawn from the material being sewn, a set of "feed dogs" is pushed up through slots in the machine surface, then dragged horizontally past the needle. The dogs are serrated to grip the material, and a "presser foot" is used to keep the material in contact with the dogs. At the end of their horizontal motion, the dogs are lowered again and returned to their original position while the needle makes its next pass through the material. While the needle is in the material, there is no feed action. Almost all household machines and the majority of industrial machines use drop feed. Differential feed is a variation of drop feed with two independent sets of dogs, one before and one after the needle. By changing their relative motions, these sets of dogs can be used to stretch or compress the material in the vicinity of the needle. This is extremely useful when sewing stretchy material, and overlock machines (heavily used for such materials) frequently have differential feed. A needle feed, used only in industrial machines, moves the material while the needle is in the material. In fact, the needle may be the primary feeding force. Some implementations of needle feed rock the axis of needle motion back and forth, while other implementations keep the axis vertical while moving it forward and back. In both cases, there is no feed action while the needle is out of the material. Needle feed is often used in conjunction with a modified drop feed, and is very common on industrial two needle machines. Most household machines do not use needle feed. A walking foot replaces the stationary presser foot with one that moves along with whatever other feed mechanisms the machine already has. As the walking foot moves, it shifts the workpiece along with it. It is most useful for sewing heavy materials where needle feed is mechanically inadequate, for spongy or cushioned materials where lifting the foot out of contact with the material helps in the feeding action, and for sewing many layers together where a drop feed will cause the lower layers to shift out of position with the upper layers. Some factory machines and a few household machines are set up with an auxiliary puller feed, which grips the material being sewn (usually from behind the needles) and pulls it with a force and reliability usually not possible with other types of feed. Puller feeds are seldom built directly into the basic sewing machine. Their action must be synchronized with the needle and feed action built into the machine to avoid damaging the machine. Pullers are also limited to straight seams, or very nearly so. Despite their additional cost and limitations, pulling feeds are very useful when making large heavy items like tents and vehicle covers. A manual feed is used primarily in freehand embroidery, quilting, and shoe repair. With manual feed, the stitch length and direction is controlled entirely by the motion of the material being sewn. Frequently some form of hoop or stabilizing material is used with fabric to keep the material under proper tension and aid in moving it around. Most household machines can be set for manual feed by disengaging the drop feed dogs. Most industrial machines can not be used for manual feed without actually removing the feed dogs. Needles Sewing machines use special needles tailored to their needs and to the character of the material being sewn. Modern sewing machines may be equipped with a needle guard. Needle guards are a safety measure that are used to help avoid injuries. Tension Tension in a sewing machine refers to the pull of the thread between the needle and the bobbin. Sewing machines have tension discs and a tension regulator. If the stitch is too saggy or too tight, the most likely cause is a tension problem. Industrial versus domestic There are mainly two types of sewing machines available: industrial and domestic. Industrial sewing machines are larger, faster, and more varied in their size, cost, appearance, and task. An industrial sewing machine can handle heavy-duty sewing jobs. Industrial machines, unlike domestic machines, perform a single dedicated task and are capable of continuous use for long periods; they have larger moving parts and larger motors rated for continuous operation. Parts for different industrial machines, such as motors, sewing feet, and bobbins may be interchangeable, but this is not always so. The motors on industrial machines, as with most of their components, lights, etc., are separate, usually mounted to the underside of the table. Domestic machines have their OEM motors mounted inside the machine. There are two different types of motor available for industrial machines: a servo motor (which uses less electricity and is silent when not in use), and the more traditional clutch motor (which is always spinning, even when not in use). A clutch motor is always running and making noises when it is connected to electricity. The constant operation ensures consistency and speed. The servo motor uses less electricity than a clutch motor. It does not make any sound unless the operator hits the pedal on the machine, but it cannot withstand the same kind of use as a clutch motor. Social impact Before sewing machines were invented women spent much of their time maintaining their family's clothing. Middle-class housewives, even with the aid of a hired seamstress, would devote several days of each month to this task. It took an experienced seamstress at least 14 hours to make a dress shirt for a man; a woman's dress took 10 hours; and a pair of summer trousers took nearly three hours. Most people except the very well-off would have only two sets of clothing: a work outfit and a Sunday outfit. Sewing machines reduced the time for making a dress shirt to an hour and 15 minutes; the time to make a dress to an hour; and the time for a pair of summer pants to 38 minutes. This reduced labor resulted in women having a diminished role in household management, and allowed more hours for their own leisure as well as the ability to seek more employment. Industrial use of sewing machines further reduced the burden placed upon housewives, moving clothing production from housewives and seamstresses to large-scale factories. The movement to large-scale factories resulted in a great increase in productivity; fewer workers could produce the same amount of clothing, reducing clothing prices significantly. As supply increased, prices also dropped. While many middle-class women enjoyed increased leisure during the Victorian era, working-class women faced intensifying demands, particularly in the clothing industry. The invention of the sewing machine, which revolutionized garment production, brought longer working hours for seamstresses, especially during peak times of the year when wealthy customers placed orders in preparation for "the Season”—the high point of the fashion season. Many women worked a minimum of eighteen hours, and sometimes up to twenty hours a day, particularly those employed by high-end London retailers. Despite the technological advancements, the industry’s seasonality left women with low wages during peak periods and no income for much of the year. Faced with such economic instability, some women turned to prostitution to survive. This moral and economic tension extended to broader societal concerns. A challenge working-class women in the textile industry often faced was the notion that factory work could lead to moral decline, fueling fears of widespread prostitution. This anxiety was reflected in literature of the time, with several novels depicting female characters who fell into prostitution after entering factory work. Often, these narratives attributed their downfall to spending wages on ‘finery’ or clothing, which was seen as fostering vanity and eventual vice. Elizabeth Gaskell’s Mary Barton (1848), set in Manchester, exemplifies these themes. From the 1840s to the 1920s, such questions about the relationship between the clothing women wore and the clothing they made generated moral panics. Works like Benjamin Disraeli's Sybil (1845), Theodore Dreiser's Sister Carrie (1900), Emile Zola's Au Bonheur des Dames (1883), and Honoré de Balzac's Splendeurs et misères des courtisanes (1847) similarly explored these themes. Adding to these discussions was the influence of William Acton, a Parisian-trained British doctor who argued that one of the primary causes of prostitution was women’s excessive love of finery, or clothing. Acton suggested that this "vanity" drove women into prostitution, although historians debate whether he meant that women engaged in prostitution to buy fashionable clothing or that admiration for the attire worn by prostitutes led them into the practice. Public debates surrounding the Contagious Disease Acts of 1864 and 1866 highlighted opposing views. For instance, Lucy Bull, a matron from the Royal Albert Hospital, rejected Acton’s moralistic interpretation and instead attributed prostitution to the poverty many women faced. Such interpretations reflect broader societal anxieties about women’s labor and morality during this period. Helen Rogers (1997) observes that by midcentury, the needlewoman had become an iconic figure in the Victorian imagination. She symbolized isolation, sexuality, single women’s work, wealthy women’s vanity, and even prostitution, blending labor and morality in complex ways. This perception of seamstresses often mirrored broader concerns about women working outside the home and the implications for societal norms. Pam Inder (2015) highlights that such concerns fueled a push for men to receive higher wages, enabling women to remain at home and reinforcing the Victorian belief that women should not work for pay outside the domestic sphere. However, seamstresses, a dominant segment of working women, remained at the center of these tensions. Their work was frequently viewed negatively, symbolizing a broader fear of women’s growing participation in the workforce and its perceived moral risks. Despite these societal ideals, the economic reality of women’s work remained stark. Louise Tilly and Joan Scott (1987) note that nearly half of women workers in England held manufacturing jobs in 1851. Crucially, 40 percent of these women worked from home, performing non-mechanized outwork or slop-work, which was cheap, ready-made clothing. Home-based production offered significant advantages to the garment industry, allowing it to rely on cheap, flexible labor. The long-standing association of needlework with women’s work reinforced this gendered division of labor, enabling industry leaders to exploit a vast network of low-cost laborers. Even as garment factories emerged in the 1850s, the industry continued to combine factory production with sweatshops and home-based labor, maintaining an exploitative and gendered system. At the same time, economic pressures frequently forced working-class families to defy legal efforts to keep children in school. Mid-19th-century laws sought to regulate child labor and prioritize education, but these initiatives often conflicted with the realities faced by rural families. Many parents sent their daughters to work for wages, often far from home, as a necessary means of survival. For working-class women, labor in the textile industry thus became both a necessity and a source of moral scrutiny, reflecting the tensions between societal ideals and economic imperatives. Many of the women who had previously been busy at home could now seek employment in factories, increasing the income for their family. This allowed families to be able to afford more sets of clothing and items than they previously could. For seamstresses, home sewing machines allowed them to produce clothing for the average person during periods when demand for fitted clothes was low, effectively increasing their earnings. When industrial sewing machines initially became popular many seamstresses, either working in factories or from home, lost their jobs as fewer workers could now produce the same output. In the long run these now unemployed skilled workers along with thousands of men and children would eventually be able to gain employment in jobs created as the clothing industry grew. The sewing machine's effects on the clothing industry resulted in major changes for other industries as well. Cotton production needed to increase in order to match the demand of the new clothing factories. As a result, cotton became planted in new areas where it had not previously been farmed. Other industries involved in the process benefited as well such as metal companies who provided parts for the machines, and shippers to move the increased amounts of goods. In addition to being important for clothing production, sewing machines also became important in the manufacturing of furniture with upholstery, curtains and towels, toys, books, and many other products.	Technology	Household appliances	null
85029	https://en.wikipedia.org/wiki/Chemical%20synthesis	Chemical synthesis	Chemical synthesis (chemical combination) is the artificial execution of chemical reactions to obtain one or several products. This occurs by physical and chemical manipulations usually involving one or more reactions. In modern laboratory uses, the process is reproducible and reliable. A chemical synthesis involves one or more compounds (known as reagents or reactants) that will experience a transformation under certain conditions. Various reaction types can be applied to formulate a desired product. This requires mixing the compounds in a reaction vessel, such as a chemical reactor or a simple round-bottom flask. Many reactions require some form of processing ("work-up") or purification procedure to isolate the final product. The amount produced by chemical synthesis is known as the reaction yield. Typically, yields are expressed as a mass in grams (in a laboratory setting) or as a percentage of the total theoretical quantity that could be produced based on the limiting reagent. A side reaction is an unwanted chemical reaction that can reduce the desired yield. The word synthesis was used first in a chemical context by the chemist Hermann Kolbe. Strategies Chemical synthesis employs various strategies to achieve efficient, precise, and molecular transformations that are more complex than simply converting a reactant A to a reaction product B directly. These strategies can be grouped into approaches for managing reaction sequences. Reaction Sequences: Multistep synthesis involves sequential chemical reactions, each requiring its own work-up to isolate intermediates before proceeding to the next stage. For example, the synthesis of paracetamol typically requires three separate reactions. Divergent synthesis starts with a common intermediate, which branches into multiple final products through distinct reaction pathways. Convergent synthesis synthesis involves the combination of multiple intermediates synthesized independently to create a complex final product. One-pot synthesis involves multiple reactions in the same vessel, allowing sequential transformations without intermediate isolation, reducing material loss, time, and the need for additional purification. Cascade reactions, a specific type of one-pot synthesis, streamline the process further by enabling consecutive transformations within a single reactant, minimizing resource consumption Catalytic Strategies: Catalysts play a vital role in chemical synthesis by accelerating reactions and enabling specific transformations. Photoredox catalysis provides enhanced control over reaction conditions by regulating the activation of small molecules and the oxidation state of metal catalysts. Biocatalysis uses enzymes as catalysts to speed up chemical reactions with high specificity under mild conditions. Reactivity Control: Chemoselectivity ensures that a specific functional group in a molecule reacts while others remain unaffected. Protecting groups temporarily mask reactive sites to enable selective reactions. Kinetic control prioritizes reaction pathways that form products quickly, often yielding less stable compounds. In contrast, thermodynamic control favors the formation of the most stable products. Advanced Planning and Techniques: Retrosynthetic analysis is a strategy used to plan complex syntheses by breaking down the target molecule into simpler precursors. Flow chemistry is a continuous reaction method where reactants are pumped through a reactor, allowing precise control over reaction conditions and scalability. This approach has been employed in the large-scale production of pharmaceuticals such as Tamoxifen. Organic synthesis Organic synthesis is a special type of chemical synthesis dealing with the synthesis of organic compounds. For the total synthesis of a complex product, multiple procedures in sequence may be required to synthesize the product of interest, needing a lot of time. A purely synthetic chemical synthesis begins with basic lab compounds. A semisynthetic process starts with natural products from plants or animals and then modifies them into new compounds. Inorganic synthesis Inorganic synthesis and organometallic synthesis are used to prepare compounds with significant non-organic content. An illustrative example is the preparation of the anti-cancer drug cisplatin from potassium tetrachloroplatinate. Green Chemistry Chemical synthesis using green chemistry promotes the design of new synthetic methods and apparatus that simplify operations and seeks environmentally benign solvents. Key principles include atom economy, which aims to incorporate all reactant atoms into the final product, and the reduction of waste and inefficiencies in chemical processes. Innovations in green chemistry, contribute to more sustainable and efficient chemical synthesis, reducing the environmental and health impacts of traditional methods. Applications Chemical synthesis plays a crucial role across various industries, enabling the development of materials, medicines, and technologies with significant real-world impacts. Catalysis: The development of catalysts is vital for numerous industrial processes, including petroleum refining, petrochemical production, and pollution control. Catalysts synthesized through chemical processes enhance the efficiency and sustainability of these operations. Medicine: Organic synthesis plays a vital role in drug discovery, allowing chemists to develop and optimize new drugs by modifying organic molecules. Additionally, the synthesis of metal complexes for medical imaging and cancer treatments is a key application of chemical synthesis, enabling advanced diagnostic and therapeutic techniques. Biopharmaceuticals: Chemical synthesis is critical in the production of biopharmaceuticals, including monoclonal antibodies and other biologics. Chemical synthesis enables the creation and modification of organic and biologically sourced compounds used in these treatments. Advanced techniques, such as DNA recombinant technology and cell fusion, rely on chemical synthesis to produce biologics tailored for specific diseases, ensuring they work effectively and target diseases precisely.	Physical sciences	Chemistry: General	null
85342	https://en.wikipedia.org/wiki/Coke%20%28fuel%29	Coke (fuel)	Coke is a grey, hard, and porous coal-based fuel with a high carbon content. It is made by heating coal or petroleum in the absence of air. Coke is an important industrial product, used mainly in iron ore smelting, but also as a fuel in stoves and forges. The unqualified term "coke" usually refers to the product derived from low-ash and low-sulphur bituminous coal by a process called coking. A similar product called petroleum coke, or pet coke, is obtained from crude petroleum in petroleum refineries. Coke may also be formed naturally by geologic processes. It is the residue of a destructive distillation process. Production Industrial coke furnaces The industrial production of coke from coal is called coking. The coal is baked in an airless kiln, a "coke furnace" or "coking oven", at temperatures as high as but usually around . This process vaporises or decomposes organic substances in the coal, driving off water and other volatile and liquid products such as coal gas and coal tar. Coke is the non-volatile residue of the decomposition, the cemented-together carbon and mineral residue of the original coal particles in the form of a hard and somewhat glassy solid. Additional byproducts of the coking are coal tar pitch, ammonia (NH3), hydrogen sulphide (H2S), pyridine, hydrogen cyanide and carbon based material. Some facilities have "by-product" coking ovens in which the volatile decomposition products are collected, purified and separated for use in other industries, as fuel or chemical feedstocks. Otherwise the volatile byproducts are burned to heat the coking ovens. This is an older method, but is still being used for new construction. Sources Bituminous coal must meet a set of criteria for use as coking coal, determined by particular coal assay techniques. These include moisture content, ash content, sulphur content, volatile content, tar, and plasticity. The goal is to achieve a blend of coal that when processed will produce a coke of appropriate strength (generally measured by coke strength after reaction), while losing an appropriate amount of mass. Other blending considerations include ensuring the coke will not swell too much during production and destroy the coke oven through excessive wall pressures. The greater the volatile matter in coal, the more by-product can be produced. It is generally considered that levels of 26–29% of volatile matter in the coal blend are good for coking purposes. Thus, different types of coal are proportionally blended to reach acceptable levels of volatility before the coking process begins. If the range of coal types is too great, the resulting coke is of widely varying strength and ash content, and is usually unsaleable, although in some cases it may be sold as an ordinary heating fuel. As coke has already lost its volatile matter, it cannot be coked again. Coking coal is different from thermal coal, but arises from the same basic coal-forming process. Coking coal has different macerals from thermal coal, i.e. different forms of the compressed and fossilized vegetative matter that compose the coal. The different macerals arise from different mixtures of the plant species, and variations of the conditions under which the coal has formed. Coking coal is graded according to its ash percentage-by-weight after burning: Steel Grade I (Ash content not exceeding 15%) Steel Grade II (Exceeding 15% but not exceeding 18%) Washery Grade I (Exceeding 18% but not exceeding 21%) Washery Grade II (Exceeding 21% but not exceeding 24%) Washery Grade III (Exceeding 24% but not exceeding 28%) Washery Grade IV (Exceeding 28% but not exceeding 35%) The "hearth" process The "hearth" process of coke-making, using lump coal, was akin to that of charcoal-burning; instead of a heap of prepared wood, covered with twigs, leaves and earth, there was a heap of coal, covered with coke dust. The hearth process continued to be used in many areas during the first half of the 19th century, but two events greatly lessened its importance. These were the invention of the hot blast in iron-smelting and the introduction of the beehive coke oven. The use of a blast of hot air, instead of cold air, in the smelting furnace was first introduced by Neilson in Scotland in 1828. The hearth process of making coke from coal is a very lengthy process. Beehive coke oven A fire brick chamber shaped like a dome is used, commonly known as a beehive oven. It is typically about wide and high. The roof has a hole for charging the coal or other kindling from the top. A discharging hole is provided in the circumference of the lower part of the wall. In a coke oven battery, a number of ovens are built in a row with common walls between neighboring ovens. A battery consisted of a great many ovens, sometimes hundreds, in a row. Coal is introduced from the top to produce an even layer of about deep. Air is supplied initially, to ignite the coal. Carbonization starts and produces volatile matter, which burns inside the partially closed side door. Carbonization proceeds from top to bottom and is completed in two to three days. The heat required for the process is supplied by the burning volatile matter, so no by-products are recovered. The exhaust gases are allowed to escape to the atmosphere. The hot coke is quenched with water, and is discharged manually through the side door. When the oven is used on a continuous basis, the walls and roof retain enough heat to initiate carbonization of the next charge. When coal was burned in a coke oven, the impurities of the coal that were not driven off as gases accumulated in the oven as slag – effectively a conglomeration of the removed impurities. Since this slag was not the desired product, it was initially just discarded. Later, however, coke oven slag was found to be useful, and has since been used as an ingredient in brick-making, mixed cement, granule-covered shingles, and even as a fertilizer. Occupational safety People can be exposed to coke oven emissions in the workplace by inhalation, skin contact, or eye contact. For the United States, the Occupational Safety and Health Administration (OSHA) has set the legal limit for coke oven emissions exposure in the workplace as 0.150 mg/m3 benzene-soluble fraction over an eight-hour workday. The US National Institute for Occupational Safety and Health (NIOSH) has set a recommended exposure limit (REL) of 0.2 mg/m3 benzene-soluble fraction over an eight-hour workday. Uses Coke can be used as a fuel and as a reducing agent in smelting iron ore in a blast furnace. The carbon monoxide produced by combustion of coke reduces iron oxide (hematite) to produce iron: Fe2O3 + 3CO -> 2Fe + 3CO2. Coke is commonly used as fuel for blacksmithing. Coke was used in Australia in the 1960s and early 1970s for house heating, and was incentivized for home use in the UK (so as to displace coal) after the 1956 Clean Air Act, which was passed in response to the Great Smog of London in 1952. Since smoke-producing constituents are driven off during the coking of coal, coke forms a desirable fuel for stoves and furnaces in which conditions are not suitable for the complete burning of bituminous coal itself. Coke may be combusted producing little or no smoke, while bituminous coal would produce much smoke. Coke was widely used as a smokeless fuel substitute for coal in domestic heating following the creation of "smokeless zones" in the United Kingdom. Highland Park distillery in Orkney roasts malted barley for use in their Scotch whisky in kilns burning a mixture of coke and peat. Coke may be used to make synthesis gas, a mixture of carbon monoxide and hydrogen. Syngas; water gas: a mixture of carbon monoxide and hydrogen, made by passing steam over red-hot coke (or any carbon-based char). Hydrocarbonate (gas) is identical, although it emerged in the late eighteenth century as an inhalation therapeutic developed by Thomas Beddoes and James Watt categorized under factitious airs Producer gas; wood gas; generator gas; synthetic gas: a mixture of carbon monoxide, hydrogen, and nitrogen, made by passing air over red-hot coke (or any carbon-based char) Coke oven gas generated from coke ovens is similar to Syngas with 60% hydrogen by volume. The hydrogen can be extracted from the coke oven gas economically for various uses (including steel production). In foundry components Finely ground bituminous coal, known in this application as sea coal, is a constituent of foundry sand. While the molten metal is in the mould, the coal burns slowly, releasing reducing gases at pressure, and so preventing the metal from penetrating the pores of the sand. It is also contained in 'mould wash', a paste or liquid with the same function applied to the mould before casting. Sea coal can be mixed with the clay lining (the "bod") used for the bottom of a cupola furnace. When heated, the coal decomposes and the bod becomes slightly friable, easing the process of breaking open holes for tapping the molten metal. Phenolic byproducts Wastewater from coking is highly toxic and carcinogenic. It contains phenolic, aromatic, heterocyclic, and polycyclic organics, and inorganics including cyanides, sulfides, ammonium and ammonia. Various methods for its treatment have been studied in recent years. The white rot fungus Phanerochaete chrysosporium can remove up to 80% of phenols from coking waste water. Properties Before bituminous coal is used as coking coal, it must meet a set of criteria determined by particular coal assay techniques. The bulk specific gravity of coke is typically around 0.77. It is highly porous. Both the chemical composition and physical properties are important to the usefulness of coke in blast furnaces. In terms of composition, low ash and sulphur content are desirable. Other important characteristics are the M10, M25, and M40 test crush indexes, which convey the strength of coke during transportation into the blast furnaces; depending on the blast furnace's size, finely crushed coke pieces must not be allowed into the furnace because they would impede the flow of gas through the charge of iron and coke. A related characteristic is the Coke Strength After Reaction (CSR) index; it represents coke's ability to withstand the violent conditions inside the blast furnace before turning into fine particles. Pieces of coke are denoted with the following terminology: "bell coke" (30 - 80 mm), "nut coke" (10 - 30 mm), "coke breeze" (< 10 mm). The water content in coke is practically zero at the end of the coking process, but it is often water quenched so that it can be transported to the blast furnaces. The porous structure of coke absorbs some water, usually 3–6% of its mass. In more modern coke plants an advanced method of coke cooling uses air quenching. Other processes The solid residue remaining from refinement of petroleum by the "cracking" process is also a form of coke. Petroleum coke has many uses besides being a fuel, such as the manufacture of dry cells and of electrolytic and welding electrodes. Gas works manufacturing syngas also produce coke as an end product, called gas house coke. Fluid coking is a process which converts heavy residual crude into lighter products such as naphtha, kerosene, heating oil, and hydrocarbon gases. The "fluid" term refers to the fact that solid coke particles behave as a fluid solid in the continuous fluid coking process versus the older batch delayed-coking process where a solid mass of coke builds up in the coke drum over time. Due to a lack of oil or high-quality coals in East Germany, scientists developed a process to turn low-quality lignite into coke called high temperature lignite coke. Alternatives to coke Scrap steel can be recycled in an electric arc furnace; and an alternative to making iron by smelting is direct reduced iron, where any carbonaceous fuel can be used to make sponge or pelletised iron. To lessen carbon dioxide emissions hydrogen can be used as the reducing agent and biomass or waste as the source of carbon. Historically, charcoal has been used as an alternative to coke in a blast furnace, with the resultant iron being known as charcoal iron. History China Many historical sources dating to the 4th century describe the production of coke in ancient China. The Chinese first used coke for heating and cooking no later than the 9th century. By the first decades of the 11th century, Chinese ironworkers in the Yellow River valley began to fuel their furnaces with coke, solving their fuel problem in that tree-sparse region. By 1078 CE, the implementation of coke as a replacement to charcoal in the production of iron in China dramatically increased the industry to 125,000 tons per year. The iron was used for the creation of tools, weapons, chains for suspension bridges, and Buddhist statues. China is the largest producer and exporter of coke today. China produces 60% of the world's coke. Concerns about air pollution have motivated technological changes in the coke industry by elimination of outdated coking technologies that are not energy-efficient. Britain In 1589, a patent was granted to Thomas Proctor and William Peterson for making iron and steel and melting lead with "earth-coal, sea-coal, turf, and peat". The patent contains a distinct allusion to the preparation of coal by "cooking". In 1590, a patent was granted to the Dean of York to "purify pit-coal and free it from its offensive smell". In 1620, a patent was granted to a company composed of William St. John and other knights, mentioning the use of coke in smelting ores and manufacturing metals. In 1627, a patent was granted to Sir John Hacket and Octavius de Strada for a method of rendering sea-coal and pit-coal as useful as charcoal for burning in houses, without offense by smell of smoke. In 1603, Hugh Plat suggested that coal might be charred in a manner analogous to the way charcoal is produced from wood. This process was not employed until 1642, when coke was used for roasting malt in Derbyshire; previously, brewers had used wood, as uncoked coal cannot be used in brewing because its sulphurous fumes would impart a foul taste to the beer. It was considered an improvement in quality, and brought about an "alteration which all England admired"—the coke process allowed for a lighter roast of the malt, leading to the creation of what by the end of the 17th century was called pale ale. In 1709, Abraham Darby I established a coke-fired blast furnace to produce cast iron. Coke's superior crushing strength allowed blast furnaces to become taller and larger. The ensuing availability of inexpensive iron was one of the factors leading to the Industrial Revolution. Before this time, iron-making used large quantities of charcoal, produced by burning wood. As the coppicing of forests became unable to meet the demand, the substitution of coke for charcoal became common in Great Britain, and coke was manufactured by burning coal in heaps on the ground so that only the outer layer burned, leaving the interior of the pile in a carbonized state. In the late 18th century, brick beehive ovens were developed, which allowed more control over the burning process. In 1768, John Wilkinson built a more practical oven for converting coal into coke. Wilkinson improved the process by building the coal heaps around a low central chimney built of loose bricks and with openings for the combustion gases to enter, resulting in a higher yield of better coke. With greater skill in the firing, covering and quenching of the heaps, yields were increased from about 33% to 65% by the middle of the 19th century. The Scottish iron industry expanded rapidly in the second quarter of the 19th century, through the adoption of the hot-blast process in its coalfields. In 1802, a battery of beehive ovens was set up near Sheffield, to coke the Silkstone coal seam for use in crucible steel melting. By 1870, there were 14,000 beehive ovens in operation on the West Durham coalfields, producing 4,000,000 long tons of coke per year. As a measure of the expansion of coke making, the requirements of the iron industry in Britain were about 1,000,000 tons per year in the early 1850s, rising to about 7,000,000 tons by 1880. Of these, about 5,000,000 tons were produced in Durham county, 1,000,000 tons in the South Wales coalfield, and 1,000,000 tons in Yorkshire and Derbyshire. In the first years of steam locomotives, coke was the normal fuel. This resulted from an early piece of environmental legislation; any proposed locomotive had to "consume its own smoke". This was not technically possible to achieve until the firebox arch came into use, but burning coke, with its low smoke emissions, was considered to meet the requirement. This rule was quietly dropped, and cheaper coal became the normal fuel, as railways gained acceptance among the public. The smoke plume produced by a travelling locomotive seems now to be a mark of a steam railway, and so preserved for posterity. So-called "gas works" produced coke by heating coal in enclosed chambers. The flammable gas that was given off was stored in gas holders, to be used domestically and industrially for cooking, heating and lighting. The gas was commonly known as "town gas" since underground networks of pipes ran through most towns. It was replaced by "natural gas" (initially from the North Sea oil and gas fields) in the decade after 1967. Other byproducts of coke production included tar and ammonia, while the coke was used instead of coal in cooking ranges and to provide heat in domestic premises before the advent of central heating. United States In the US, the first use of coke in an iron furnace occurred around 1817 at Isaac Meason's Plumsock puddling furnace and rolling mill in Fayette County, Pennsylvania. In the late 19th century, the coalfields of western Pennsylvania provided a rich source of raw material for coking. In 1885, the Rochester and Pittsburgh Coal and Iron Company constructed the world's longest string of coke ovens in Walston, Pennsylvania, with 475 ovens over a length of 2 km (1.25 miles). Their output reached 22,000 tons per month. The Minersville Coke Ovens in Huntingdon County, Pennsylvania, were listed on the National Register of Historic Places in 1991. Between 1870 and 1905, the number of beehive ovens in the US increased from approximately 200 to nearly 31,000, which produced nearly 18,000,000 tons of coke in the Pittsburgh area alone. One observer boasted that if loaded into a train, "the year's production would make up a train so long that the engine in front of it would go to San Francisco and come back to Connellsville before the caboose had gotten started out of the Connellsville yards!" The number of beehive ovens in Pittsburgh peaked in 1910 at almost 48,000. Although it made a top-quality fuel, coking poisoned the surrounding landscape. After 1900, the serious environmental damage of beehive coking attracted national notice, although the damage had plagued the district for decades. "The smoke and gas from some ovens destroy all vegetation around the small mining communities", noted W. J. Lauck of the U.S. Immigration Commission in 1911. Passing through the region on train, University of Wisconsin president Charles Van Hise saw "long rows of beehive ovens from which flame is bursting and dense clouds of smoke issuing, making the sky dark. By night, the scene is rendered indescribably vivid by these numerous burning pits. The beehive ovens make the entire region of coke manufacture one of dulled sky: cheerless and unhealthful." In 2024, an investigation into 17 coke burning facilities in the US could be responsible for an estimated 892 premature deaths every year, as well as increased asthma symptoms and other health impacts for residents.	Technology	Fuel	null
85385	https://en.wikipedia.org/wiki/Atherosclerosis	Atherosclerosis	Atherosclerosis is a pattern of the disease arteriosclerosis, characterized by development of abnormalities called lesions in walls of arteries. This is a chronic inflammatory disease involving many different cell types and driven by elevated levels of cholesterol in the blood. These lesions may lead to narrowing of the arterial walls due to buildup of atheromatous plaques. At the onset there are usually no symptoms, but if they develop, symptoms generally begin around middle age. In severe cases, it can result in coronary artery disease, stroke, peripheral artery disease, or kidney disorders, depending on which body part(s) the affected arteries are located in the body. The exact cause of atherosclerosis is unknown and is proposed to be multifactorial. Risk factors include abnormal cholesterol levels, elevated levels of inflammatory biomarkers, high blood pressure, diabetes, smoking (both active and passive smoking), obesity, genetic factors, family history, lifestyle habits, and an unhealthy diet. Plaque is made up of fat, cholesterol, immune cells, calcium, and other substances found in the blood. The narrowing of arteries limits the flow of oxygen-rich blood to parts of the body. Diagnosis is based upon a physical exam, electrocardiogram, and exercise stress test, among others. Prevention guidelines include eating a healthy diet, exercising, not smoking, and maintaining normal body weight. Treatment of established atherosclerotic disease may include medications to lower cholesterol such as statins, blood pressure medication, and anticoagulant therapies to reduce the risk of blood clot formation. As the disease state progresses more invasive strategies are applied such as percutaneous coronary intervention, coronary artery bypass graft, or carotid endarterectomy. Genetic factors are also strongly implicated in the disease process; it is unlikely to be entirely based on lifestyle choices. Atherosclerosis generally starts when a person is young and worsens with age. Women are 78% at higher risk level than men Almost all people are affected to some degree by the age of 65. It is the number one cause of death and disability in developed countries. Though it was first described in 1575, there is evidence suggesting that this disease state is genetically inherent in the broader human population, with its origins tracing back to genetic mutations that may have occurred more than two million years ago during the evolution of hominin ancestors of modern human beings. Signs and symptoms Atherosclerosis is typically asymptomatic for decades because the arteries enlarge at all plaque locations, thus there is no effect on blood flow. Even most plaque ruptures do not produce symptoms until enough narrowing or closure of an artery, due to clots, occurs. Signs and symptoms only occur after severe narrowing or closure impedes blood flow to different organs enough to induce symptoms. Most of the time, patients realize that they have the disease only when they experience other cardiovascular disorders such as stroke or heart attack. These symptoms, however, still vary depending on which artery or organ is affected. Early atherosclerotic processes likely begin in childhood. Fibrous and gelatinous lesions have been observed in the coronary arteries of children. Fatty streaks have been observed in the coronary arteries of juveniles. While coronary artery disease is more prevalent in men than women, atherosclerosis of the cerebral arteries and strokes equally affect both sexes. Marked narrowing in the coronary arteries, which are responsible for bringing oxygenated blood to the heart, can produce symptoms such as chest pain of angina and shortness of breath, sweating, nausea, dizziness or lightheadedness, breathlessness or palpitations. Abnormal heart rhythms called arrhythmias—the heart beating either too slowly or too quickly—are another consequence of ischemia. Carotid arteries supply blood to the brain and neck. Marked narrowing of the carotid arteries can present with symptoms such as a feeling of weakness; being unable to think straight; difficulty speaking; dizziness; difficulty in walking or standing up straight; blurred vision; numbness of the face, arms and legs; severe headache; and loss of consciousness. These symptoms are also related to stroke (death of brain cells). Stroke is caused by marked narrowing or closure of arteries going to the brain; lack of adequate blood supply leads to the death of the cells of the affected tissue. Peripheral arteries, which supply blood to the legs, arms, and pelvis, also experience marked narrowing due to plaque rupture and clots. Symptoms of the narrowing are pain and numbness within the arms or legs. Another significant location for plaque formation is the renal arteries, which supply blood to the kidneys. Plaque occurrence and accumulation lead to decreased kidney blood flow and chronic kidney disease, which, like in all other areas, is typically asymptomatic until late stages. In 2004, US data indicated that in ~66% of men and ~47% of women, the first symptom of atherosclerotic cardiovascular disease was a heart attack or sudden cardiac death (defined as death within one hour of onset of the symptom). Case studies have included autopsies of U.S. soldiers killed in World War II and the Korean War. A much-cited report involved the autopsies of 300 U.S. soldiers killed in Korea. Although the average age of the men was 22.1 years, 77.3 percent had "gross evidence of coronary arteriosclerosis". Risk factors The atherosclerotic process is not well understood. Atherosclerosis is associated with inflammatory processes in the endothelial cells of the vessel wall associated with retained low-density lipoprotein (LDL) particles. This retention may be a cause, an effect, or both, of the underlying inflammatory process. The presence of the plaque induces the muscle cells of the blood vessel to stretch, compensating for the additional bulk. The endothelial lining then thickens, increasing the separation between the plaque and lumen. The thickening somewhat offsets the narrowing caused by the growth of the plaque, but moreover, it causes the wall to stiffen and become less compliant to stretching with each heartbeat. Modifiable Western pattern diet Abdominal obesity Insulin resistance Diabetes Dyslipidemia Hypertension Trans fat Tobacco smoking Bacterial infections HIV/AIDS Nonmodifiable South Asian descent Advanced age Genetic abnormalities Family history Coronary anatomy and branch pattern Lesser or uncertain Thrombophilia Saturated fat Excessive carbohydrates Elevated triglycerides Systemic inflammation Hyperinsulinemia Sleep deprivation Air pollution Sedentary lifestyle Arsenic poisoning Alcohol Chronic stress Hypothyroidism Periodontal disease Dietary The relation between dietary fat and atherosclerosis is controversial. The USDA, in its food pyramid, promotes a diet of about 64% carbohydrates from total calories. The American Heart Association, the American Diabetes Association, and the National Cholesterol Education Program make similar recommendations. In contrast, Prof Walter Willett (Harvard School of Public Health, PI of the second Nurses' Health Study) recommends much higher levels of fat, especially of monounsaturated and polyunsaturated fat. These dietary recommendations reach a consensus, though, against consumption of trans fats. The role of eating oxidized fats (rancid fats) in humans is not clear. Rabbits fed rancid fats develop atherosclerosis faster. Rats fed DHA-containing oils experienced marked disruptions to their antioxidant systems, and accumulated significant amounts of phospholipid hydroperoxide in their blood, livers and kidneys. Rabbits fed atherogenic diets containing various oils were found to undergo the most oxidative susceptibility of LDL via polyunsaturated oils. In another study, rabbits fed heated soybean oil "grossly induced atherosclerosis and marked liver damage were histologically and clinically demonstrated." However, Fred Kummerow claims that it is not dietary cholesterol, but oxysterols, or oxidized cholesterols, from fried foods and smoking, that are the culprit. Rancid fats and oils taste very unpleasant in even small amounts, so people avoid eating them. It is very difficult to measure or estimate the actual human consumption of these substances. Highly unsaturated omega-3 rich oils such as fish oil, when being sold in pill form, can hide the taste of oxidized or rancid fat that might be present. In the US, the health food industry's dietary supplements are self-regulated and outside of FDA regulations. To protect unsaturated fats from oxidation, it is best to keep them cool and in oxygen-free environments. Pathophysiology Atherogenesis is the developmental process of atheromatous plaques. It is characterized by a remodeling of arteries leading to subendothelial accumulation of fatty substances called plaques. The buildup of an atheromatous plaque is a slow process, developed over several years through a complex series of cellular events occurring within the arterial wall and in response to several local vascular circulating factors. One recent hypothesis suggests that, for unknown reasons, leukocytes, such as monocytes or basophils, begin to attack the endothelium of the artery lumen in cardiac muscle. The ensuing inflammation leads to the formation of atheromatous plaques in the arterial tunica intima, a region of the vessel wall located between the endothelium and the tunica media. Chronic inflammation within the arterial wall, driven by immune cells like macrophages, accelerates atherosclerotic plaque instability by promoting collagen breakdown and thinning the fibrous cap, which increases the likelihood of rupture and thrombosis. The bulk of these lesions is made of excess fat, collagen, and elastin. At first, as the plaques grow, only wall thickening occurs without narrowing. Stenosis is a late event, which may never occur and is often the result of repeated plaque rupture and healing responses, not just the atherosclerotic process. Autopsy studies have shown that the prevalence of coronary artery atherosclerosis in males from the United States, with an average age of 22.1 years old, who died in war, ranges from 45% to 77.3%. Cellular Early atherogenesis is characterized by the adherence of blood circulating monocytes (a type of white blood cell) to the vascular bed lining, the endothelium, then by their migration to the sub-endothelial space, and further activation into monocyte-derived macrophages. The primary documented driver of this process is oxidized lipoprotein particles within the wall, beneath the endothelial cells, though upper normal or elevated concentrations of blood glucose also plays a major role and not all factors are fully understood. Fatty streaks may appear and disappear. Low-density lipoprotein (LDL) particles in blood plasma invade the endothelium and become oxidized, creating risk of cardiovascular disease. A complex set of biochemical reactions regulates the oxidation of LDL, involving enzymes (such as Lp-LpA2) and free radicals in the endothelium. Initial damage to the endothelium results in an inflammatory response. Monocytes enter the artery wall from the bloodstream, with platelets adhering to the area of insult. This may be promoted by redox signaling induction of factors such as VCAM-1, which recruits circulating monocytes, and M-CSF, which is selectively required for the differentiation of monocytes to macrophages. The monocytes differentiate into macrophages, which proliferate locally, ingest oxidized LDL, slowly turning into large "foam cells" – so-called because of their changed appearance resulting from the numerous internal cytoplasmic vesicles and resulting high lipid content. Under the microscope, the lesion now appears as a fatty streak. Foam cells eventually die and further propagate the inflammatory process. In addition to these cellular activities, there is also smooth muscle proliferation and migration from the tunica media into the intima in response to cytokines secreted by damaged endothelial cells. This causes the formation of a fibrous capsule covering the fatty streak. Intact endothelium can prevent this smooth muscle proliferation by releasing nitric oxide. Calcification and lipids Calcification forms among vascular smooth muscle cells of the surrounding muscular layer, specifically in the muscle cells adjacent to atheromas and on the surface of atheroma plaques and tissue. In time, as cells die, this leads to extracellular calcium deposits between the muscular wall and outer portion of the atheromatous plaques. With the atheromatous plaque interfering with the regulation of calcium deposition, it accumulates and crystallizes. A similar form of intramural calcification, presenting the picture of an early phase of arteriosclerosis, appears to be induced by many drugs that have an antiproliferative mechanism of action (Rainer Liedtke 2008). Cholesterol is delivered into the vessel wall by cholesterol-containing low-density lipoprotein (LDL) particles. To attract and stimulate macrophages, the cholesterol must be released from the LDL particles and oxidized, a key step in the ongoing inflammatory process. The process is worsened if it is insufficient high-density lipoprotein (HDL), the lipoprotein particle that removes cholesterol from tissues and carries it back to the liver. The foam cells and platelets encourage the migration and proliferation of smooth muscle cells, which in turn ingest lipids, become replaced by collagen, and transform into foam cells themselves. A protective fibrous cap normally forms between the fatty deposits and the artery lining (the intima). These capped fatty deposits (now called 'atheromas') produce enzymes that cause the artery to enlarge over time. As long as the artery enlarges sufficiently to compensate for the extra thickness of the atheroma, then no narrowing ("stenosis") of the opening ("lumen") occurs. The artery expands with an egg-shaped cross-section, still with a circular opening. If the enlargement is beyond proportion to the atheroma thickness, then an aneurysm is created. Visible features Although arteries are not typically studied microscopically, two plaque types can be distinguished: The fibro-lipid (fibro-fatty) plaque is characterized by an accumulation of lipid-laden cells underneath the intima of the arteries, typically without narrowing the lumen due to compensatory expansion of the bounding muscular layer of the artery wall. Beneath the endothelium, there is a "fibrous cap" covering the atheromatous "core" of the plaque. The core consists of lipid-laden cells (macrophages and smooth muscle cells) with elevated tissue cholesterol and cholesterol ester content, fibrin, proteoglycans, collagen, elastin, and cellular debris. In advanced plaques, the central core of the plaque usually contains extracellular cholesterol deposits (released from dead cells), which form areas of cholesterol crystals with empty, needle-like clefts. At the periphery of the plaque are younger "foamy" cells and capillaries. These plaques usually produce the most damage to the individual when they rupture. Cholesterol crystals may also play a role. The fibrous plaque is also localized under the intima, within the arterial wall resulting in thickening and expansion of the wall and, sometimes, spotty localized narrowing of the lumen with some atrophy of the muscular layer. The fibrous plaque contains collagen fibers (eosinophilic), precipitates of calcium (hematoxylinophilic), and rarely, lipid-laden cells. In effect, the muscular portion of the artery wall forms small aneurysms just large enough to hold the atheroma that are present. The muscular portion of artery walls usually remains strong, even after they have been remodeled to compensate for the atheromatous plaques. However, atheromas within the vessel wall are soft and fragile with little elasticity. Arteries constantly expand and contract with each heartbeat, i.e., the pulse. In addition, the calcification deposits between the outer portion of the atheroma and the muscular wall, as they progress, lead to a loss of elasticity and stiffening of the artery as a whole. The calcification deposits, after they have become sufficiently advanced, are partially visible on coronary artery computed tomography or electron beam tomography (EBT) as rings of increased radiographic density, forming halos around the outer edges of the atheromatous plaques, within the artery wall. On CT, >130 units on the Hounsfield scale (some argue for 90 units) has been the radiographic density usually accepted as clearly representing tissue calcification within arteries. These deposits demonstrate unequivocal evidence of the disease, relatively advanced, even though the lumen of the artery is often still normal by angiography. Rupture and stenosis Although the disease process tends to be slowly progressive over decades, it usually remains asymptomatic until an atheroma ulcerates, which leads to immediate blood clotting at the site of the atheroma ulcer. This triggers a cascade of events that leads to clot enlargement, which may quickly obstruct blood flow. A complete blockage leads to ischemia of the myocardial (heart) muscle and damage. This process is the myocardial infarction or "heart attack". If the heart attack is not fatal, fibrous organization of the clot within the lumen ensues, covering the rupture but also producing stenosis or closure of the lumen, or over time and after repeated ruptures, resulting in a persistent, usually localized stenosis or blockage of the artery lumen. Stenoses can be slowly progressive, whereas plaque ulceration is a sudden event that occurs specifically in atheromas with thinner/weaker fibrous caps that have become "unstable". Repeated plaque ruptures, ones not resulting in total lumen closure, combined with the clot patch over the rupture and healing response to stabilize the clot is the process that produces most stenoses over time. The stenotic areas often become more stable despite increased flow velocities at these narrowings. Most major blood-flow-stopping events occur at large plaques, which, before their rupture, produced little if any stenosis. From clinical trials, 20% is the average stenosis at plaques that subsequently rupture with resulting complete artery closure. Most severe clinical events do not occur at plaques that produce high-grade stenosis. From clinical trials, only 14% of heart attacks occur from artery closure at plaques producing a 75% or greater stenosis before the vessel closing. If the fibrous cap separating a soft atheroma from the bloodstream within the artery ruptures, tissue fragments are exposed and released. These tissue fragments are very clot-promoting, containing collagen and tissue factor; they activate platelets and activate the system of coagulation. The result is the formation of a thrombus (blood clot) overlying the atheroma, which obstructs blood flow acutely. With the obstruction of blood flow, downstream tissues are starved of oxygen and nutrients. If this is the myocardium (heart muscle) angina (cardiac chest pain) or myocardial infarction (heart attack) develops. Accelerated growth of plaques The distribution of atherosclerotic plaques in a part of arterial endothelium is inhomogeneous. The multiple and focal development of atherosclerotic changes is similar to that in the development of amyloid plaques in the brain and age spots on the skin. Misrepair-accumulation aging theory suggests that misrepair mechanisms play an important role in the focal development of atherosclerosis. The development of a plaque is a result of the repair of the injured endothelium. Because of the infusion of lipids into the sub-endothelium, the repair has to end up with altered remodeling of the local endothelium. This is the manifestation of a misrepair. This altered remodeling increases the susceptibility of the local endothelium to damage and reduces its repair efficiency. Consequently, this part of endothelium has an increased risk of being injured and improperly repaired. Thus, the accumulation of misrepairs of endothelium is focalized and self-accelerating. In this way, the growth of a plaque is also self-accelerating. Within a part of the arterial wall, the oldest plaque is always the biggest and is the most dangerous one to cause blockage of a local artery. Components The plaque is divided into three distinct components: The atheroma ("lump of gruel", ), which is the nodular accumulation of a soft, flaky, yellowish material at the center of large plaques, composed of macrophages nearest the lumen of the artery Underlying areas of cholesterol crystals Calcification at the outer base of older or more advanced lesions. Atherosclerotic lesions, or atherosclerotic plaques, are separated into two broad categories: Stable and unstable (also called vulnerable). The pathobiology of atherosclerotic lesions is very complicated, but generally, stable atherosclerotic plaques, which tend to be asymptomatic, are rich in extracellular matrix and smooth muscle cells. On the other hand, unstable plaques are rich in macrophages and foam cells, and the extracellular matrix separating the lesion from the arterial lumen (also known as the fibrous cap) is usually weak and prone to rupture. Ruptures of the fibrous cap expose thrombogenic material, such as collagen, to the circulation and eventually induce thrombus formation in the lumen. Upon formation, intraluminal thrombi can occlude arteries outright (e.g., coronary occlusion), but more often they detach, move into the circulation, and eventually occlude smaller downstream branches causing thromboembolism. Apart from thromboembolism, chronically expanding atherosclerotic lesions can cause complete closure of the lumen. Chronically expanding lesions are often asymptomatic until the lumen stenosis is so severe (usually over 80%) that blood supply to downstream tissue(s) is insufficient, resulting in ischemia. These complications of advanced atherosclerosis are chronic, slowly progressive, and cumulative. Most commonly, soft plaque suddenly ruptures (see vulnerable plaque), causing the formation of a thrombus that will rapidly slow or stop blood flow, leading to the death of the tissues fed by the artery in approximately five minutes. This event is called an infarction. Diagnosis Areas of severe narrowing, stenosis, detectable by angiography, and to a lesser extent "stress testing" have long been the focus of human diagnostic techniques for cardiovascular disease, in general. However, these methods focus on detecting only severe narrowing, not the underlying atherosclerosis disease. As demonstrated by human clinical studies, most severe events occur in locations with heavy plaque, yet little or no lumen narrowing present before debilitating events suddenly occur. Plaque rupture can lead to artery lumen occlusion within seconds to minutes, potential permanent debility, and sometimes sudden death. Plaques that have ruptured are called complicated lesions. The extracellular matrix of the lesion breaks, usually at the shoulder of the fibrous cap that separates the lesion from the arterial lumen, where the exposed thrombogenic components of the plaque, mainly collagen, will trigger thrombus formation. The thrombus then travels downstream to other blood vessels, where the blood clot may partially or completely block blood flow. If the blood flow is completely blocked, cell deaths occur due to the lack of oxygen supply to nearby cells, resulting in necrosis. The narrowing or obstruction of blood flow can occur in any artery within the body. Obstruction of arteries supplying the heart muscle results in a heart attack, while the obstruction of arteries supplying the brain results in an ischaemic stroke. Lumen stenosis that is greater than 75% was considered the hallmark of clinically significant disease in the past because recurring episodes of angina and abnormalities in stress tests are only detectable at that particular severity of stenosis. However, clinical trials have shown that only about 14% of clinically debilitating events occur at sites with more than 75% stenosis. Most cardiovascular events that involve sudden rupture of the atheroma plaque do not display any evident luminal narrowing. Thus, greater attention has been focused on "vulnerable plaque" from the late 1990s onwards. Besides the traditional diagnostic methods such as angiography and stress-testing, other detection techniques have been developed in the past decades for earlier detection of atherosclerotic disease. Some of the detection approaches include anatomical detection and physiologic measurement. Examples of anatomical detection methods include coronary calcium scoring by CT, carotid IMT (intimal media thickness) measurement by ultrasound, and intravascular imaging techniques, such as intravascular ultrasound (IVUS), and intravascular optical coherence tomography (OCT), allowing direct visualization of atherosclerotic plaques. Other methods include blood measurements, e.g., lipoprotein subclass analysis, HbA1c, hs-CRP, and homocysteine. Both anatomic and physiologic methods allow early detection before symptoms show up, disease staging, and tracking of disease progression. In recent years, developments in nuclear imaging techniques such as PET and SPECT have provided non-invasive ways of estimating the severity of atherosclerotic plaques. Prevention Up to 90% of cardiovascular disease may be preventable if established risk factors are avoided. Medical management of atherosclerosis first involves modification to risk factors–for example, via smoking cessation and diet restrictions. Prevention is generally by eating a healthy diet, exercising, not smoking, and maintaining a normal weight. Diet Changes in diet may help prevent the development of atherosclerosis. Tentative evidence suggests that a diet containing dairy products has no effect on or decreases the risk of cardiovascular disease. A diet high in fruits and vegetables decreases the risk of cardiovascular disease and death. Evidence suggests that the Mediterranean diet may improve cardiovascular results. There is also evidence that a Mediterranean diet may be better than a low-fat diet in bringing about long-term changes to cardiovascular risk factors (e.g., lower cholesterol level and blood pressure). A 2024 review highlighted that bioactive compounds found in Mediterranean diet components (such as olive, grape, garlic, rosemary, and saffron) exhibit properties that may contribute to cardiovascular health and atherosclerosis prevention. Exercise A controlled exercise program combats atherosclerosis by improving the circulation and functionality of the vessels. Exercise is also used to manage weight in patients who are obese, lower blood pressure, and decrease cholesterol. Often lifestyle modification is combined with medication therapy. For example, statins help to lower cholesterol. Antiplatelet medications like aspirin help to prevent clots, and a variety of antihypertensive medications are routinely used to control blood pressure. If the combined efforts of risk factor modification and medication therapy are not sufficient to control symptoms or fight imminent threats of ischemic events, a physician may resort to interventional or surgical procedures to correct the obstruction. Treatment Treatment of established disease may include medications to lower cholesterol such as statins, blood pressure medication, or medications that decrease clotting, such as aspirin. Many procedures may also be carried out such as percutaneous coronary intervention, coronary artery bypass graft, or carotid endarterectomy. Medical treatments often focus on alleviating symptoms. However, measures that focus on decreasing underlying atherosclerosis—as opposed to simply treating symptoms—are more effective. Non-pharmaceutical means are usually the first method of treatment, such as stopping smoking and practicing regular exercise. If these methods do not work, medicines are usually the next step in treating cardiovascular diseases and, with improvements, have increasingly become the most effective method over the long term. The key to the more effective approaches is to combine different treatment strategies. In addition, for those approaches, such as lipoprotein transport behaviors, which have been shown to produce the most success, adopting more aggressive combination treatment strategies taken daily and indefinitely has generally produced better results, both before and especially after people are symptomatic. Statins Statin medications are widely prescribed for treating atherosclerosis. They have shown benefit in reducing cardiovascular disease and mortality in those with high cholesterol with few side effects. Secondary prevention therapy, which includes high-intensity statins and aspirin, is recommended by multi-society guidelines for all patients with a history of ASCVD (atherosclerotic cardiovascular disease) to prevent the recurrence of coronary artery disease, ischemic stroke, or peripheral arterial disease. However, prescription of and adherence to these guideline-concordant therapies is lacking, particularly among young patients and women. Statins work by inhibiting HMG-CoA (hydroxymethylglutaryl-coenzyme A) reductase, a hepatic rate-limiting enzyme in cholesterol's biochemical production pathway. Inhibiting this rate-limiting enzyme reduces the body's ability to produce as much cholesterol endogenously, thereby reducing the level of LDL-cholesterol in the blood. This reduced endogenous cholesterol production triggers the body to then pull cholesterol from other cellular sources, enhancing serum HDL-cholesterol. These data are primarily in middle-aged men and the conclusions are less clear for women and people over the age of 70. Surgery When atherosclerosis has become severe and caused irreversible ischemia, such as tissue loss in the case of peripheral artery disease, surgery may be indicated. Vascular bypass surgery can re-establish flow around the diseased segment of the artery, and angioplasty with or without stenting can reopen narrowed arteries and improve blood flow. Coronary artery bypass grafting without manipulation of the ascending aorta has demonstrated reduced rates of postoperative stroke and mortality compared to traditional on-pump coronary revascularization. Other There is evidence that some anticoagulants, particularly warfarin, which inhibit clot formation by interfering with Vitamin K metabolism, may promote arterial calcification in the long term despite reducing clot formation in the short term. Also, small molecules such as 3-hydroxybenzaldehyde and protocatechuic aldehyde have shown vasculoprotective effects to reduce risk of atherosclerosis. Epidemiology Cardiovascular disease, which is predominantly the clinical manifestation of atherosclerosis, is one of the leading causes of death worldwide. Almost all children older than age 10 in developed countries have aortic fatty streaks, with coronary fatty streaks beginning in adolescence. In 1953, a study was published that examined the results of 300 autopsies performed on U.S. soldiers who had died in the Korean War. Despite the average age of the soldiers being just 22 years old, 77% of them had visible signs of coronary atherosclerosis. This study showed that heart disease could affect people at a younger age and was not just a problem for older individuals. In 1992, a report showed that microscopic fatty streaks were seen in the left anterior descending artery in over 50% of children aged 10–14 and 8% had even more advanced lesions with more accumulations of extracellular lipid. A 2005 report of a study done between 1985 and 1995 found that around 87% of aortas and 30% of coronary arteries in the age group 5–14 years had fatty streaks which increased with age. Etymology The following terms are similar, yet distinct, in both spelling and meaning, and can be easily confused: arteriosclerosis, arteriolosclerosis, and atherosclerosis. Arteriosclerosis is a general term describing any hardening (and loss of elasticity) of medium or large arteries (); arteriolosclerosis is any hardening (and loss of elasticity) of arterioles (small arteries); atherosclerosis is a hardening of an artery specifically due to an atheromatous plaque (). The term atherogenic is used for substances or processes that cause the formation of atheroma. Economics In 2011, coronary atherosclerosis was one of the top ten most expensive conditions seen during inpatient hospitalizations in the US, with aggregate inpatient hospital costs of $10.4 billion. Research Lipids An indication of the role of high-density lipoprotein (HDL) on atherosclerosis has been with the rare Apo-A1 Milano human genetic variant of this HDL protein. A small short-term trial using bacterial synthesized human Apo-A1 Milano HDL in people with unstable angina produced a fairly dramatic reduction in measured coronary plaque volume in only six weeks vs. the usual increase in plaque volume in those randomized to placebo. The trial was published in JAMA in early 2006. Ongoing work starting in the 1990s may lead to human clinical trials—probably by about 2008. These may use synthesized Apo-A1 Milano HDL directly, or they may use gene-transfer methods to pass the ability to synthesize the Apo-A1 Milano HDLipoprotein. Methods to increase HDL particle concentrations, which in some animal studies largely reverses and removes atheromas, are being developed and researched. However, increasing HDL by any means is not necessarily helpful. For example, the drug torcetrapib is the most effective agent currently known for raising HDL (by up to 60%). However, in clinical trials, it also raised deaths by 60%. All studies regarding this drug were halted in December 2006. The actions of macrophages drive atherosclerotic plaque progression. Immunomodulation of atherosclerosis is the term for techniques that modulate immune system function to suppress this macrophage action. Involvement of lipid peroxidation chain reaction in atherogenesis triggered research on the protective role of the heavy isotope (deuterated) polyunsaturated fatty acids (D-PUFAs) that are less prone to oxidation than ordinary PUFAs (H-PUFAs). PUFAs are essential nutrients – they are involved in metabolism in that very form as they are consumed with food. In transgenic mice, that are a model for human-like lipoprotein metabolism, adding D-PUFAs to diet indeed reduced body weight gain, improved cholesterol handling and reduced atherosclerotic damage to the aorta. miRNA MicroRNAs (miRNAs) have complementary sequences in the 3' UTR and 5' UTR of target mRNAs of protein-coding genes, and cause mRNA cleavage or repression of translational machinery. In diseased vascular vessels, miRNAs are dysregulated and highly expressed. miR-33 is found in cardiovascular diseases. It is involved in atherosclerotic initiation and progression including lipid metabolism, insulin signaling and glucose homeostatis, cell type progression and proliferation, and myeloid cell differentiation. It was found in rodents that the inhibition of miR-33 will raise HDL-C levels and the expression of miR-33 is down-regulated in humans with atherosclerotic plaques. miR-33a and miR-33b are located on intron 16 of human sterol regulatory element-binding protein 2 (SREBP2) gene on chromosome 22 and intron 17 of SREBP1 gene on chromosome 17. miR-33a/b regulates cholesterol/lipid homeostasis by binding in the 3'UTRs of genes involved in cholesterol transport such as ATP binding cassette (ABC) transporters and enhance or represses its expression. Studies have shown that ABCA1 mediates cholesterol transport from peripheral tissues to Apolipoprotein-1. It is also important in the reverse cholesterol transport pathway, where cholesterol is delivered from peripheral tissue to the liver, where it can be excreted into bile or converted to bile acids before excretion. Therefore, ABCA1 prevents cholesterol accumulation in macrophages. By enhancing miR-33 function, the level of ABCA1 is decreased, leading to decreased cellular cholesterol efflux to apoA-1. On the other hand, by inhibiting miR-33 function, the level of ABCA1 is increased and increases the cholesterol efflux to apoA-1. Suppression of miR-33 will lead to less cellular cholesterol and higher plasma HDL level through the regulation of ABCA1 expression. The sugar, cyclodextrin, removed cholesterol that had built up in the arteries of mice fed a high-fat diet. DNA damage Aging is the most important risk factor for cardiovascular problems. The causative basis by which aging mediates its impact, independently of other recognized risk factors, remains to be determined. Evidence has been reviewed for a key role of DNA damage in vascular aging. 8-oxoG, a common type of oxidative damage in DNA, is found to accumulate in plaque vascular smooth muscle cells, macrophages and endothelial cells, thus linking DNA damage to plaque formation. DNA strand breaks also increased in atherosclerotic plaques. Werner syndrome (WS) is a premature aging condition in humans. WS is caused by a genetic defect in a RecQ helicase that is employed in several repair processes that remove damages from DNA. WS patients develop a considerable burden of atherosclerotic plaques in their coronary arteries and aorta: calcification of the aortic valve is also frequently observed. These findings link excessive unrepaired DNA damage to premature aging and early atherosclerotic plaque development (see DNA damage theory of aging). Microorganisms The microbiota – all the microorganisms in the body, can contribute to atherosclerosis in many ways: modulation of the immune system, changes in metabolism, processing of nutrients and production of certain metabolites that can get into blood circulation. One such metabolite, produced by gut bacteria, is trimethylamine N-oxide (TMAO). Its levels have been associated with atherosclerosis in human studies and animal research suggests that may be a causal relation. An association between the bacterial genes encoding trimethylamine lyases — the enzymes involved in TMAO generation — and atherosclerosis has been noted. Vascular smooth muscle cells Vascular smooth muscle cells play a key role in atherogenesis and were historically considered to be beneficial for plaque stability by forming a protective fibrous cap and synthesizing strength-giving extracellular matrix components. However, in addition to the fibrous cap, vascular smooth muscle cells also give rise to many of the cell types found within the plaque core and can modulate their phenotype to both promote and reduce plaque stability. Vascular smooth muscle cells exhibit pronounced plasticity within atherosclerotic plaque and can modify their gene expression profile to resemble various other cell types, including macrophages, myofibroblasts, mesenchymal stem cells and osteochondrocytes. Importantly, genetic lineage-tracing experiments have unequivocally shown that 40-90% of plaque-resident cells are vascular smooth muscle cell-derived, therefore, it is important to research the role of vascular smooth muscle cells in atherosclerosis to identify new therapeutic targets.	Biology and health sciences	Specific diseases	Health
85411	https://en.wikipedia.org/wiki/PH%20indicator	PH indicator	A pH indicator is a halochromic chemical compound added in small amounts to a solution so the pH (acidity or basicity) of the solution can be determined visually or spectroscopically by changes in absorption and/or emission properties. Hence, a pH indicator is a chemical detector for hydronium ions (H3O+) or hydrogen ions (H+) in the Arrhenius model. Normally, the indicator causes the color of the solution to change depending on the pH. Indicators can also show change in other physical properties; for example, olfactory indicators show change in their odor. The pH value of a neutral solution is 7.0 at 25°C (standard laboratory conditions). Solutions with a pH value below 7.0 are considered acidic and solutions with pH value above 7.0 are basic. Since most naturally occurring organic compounds are weak electrolytes, such as carboxylic acids and amines, pH indicators find many applications in biology and analytical chemistry. Moreover, pH indicators form one of the three main types of indicator compounds used in chemical analysis. For the quantitative analysis of metal cations, the use of complexometric indicators is preferred, whereas the third compound class, the redox indicators, are used in redox titrations (titrations involving one or more redox reactions as the basis of chemical analysis). Theory In and of themselves, pH indicators are usually weak acids or weak bases. The general reaction scheme of acidic pH indicators in aqueous solutions can be formulated as: HInd(aq) + (l) (aq) + (aq) where, "HInd" is the acidic form and "Ind−" is the conjugate base of the indicator. Vice versa for basic pH indicators in aqueous solutions: IndOH(aq) + (l) (l) + (aq) + (aq) where "IndOH" stands for the basic form and "Ind+" for the conjugate acid of the indicator. The ratio of concentration of conjugate acid/base to concentration of the acidic/basic indicator determines the pH (or pOH) of the solution and connects the color to the pH (or pOH) value. For pH indicators that are weak electrolytes, the Henderson–Hasselbalch equation can be written as: pH = pKa + log10 pOH = pKb + log10 The equations, derived from the acidity constant and basicity constant, states that when pH equals the pKa or pKb value of the indicator, both species are present in a 1:1 ratio. If pH is above the pKa or pKb value, the concentration of the conjugate base is greater than the concentration of the acid, and the color associated with the conjugate base dominates. If pH is below the pKa or pKb value, the converse is true. Usually, the color change is not instantaneous at the pKa or pKb value, but a pH range exists where a mixture of colors is present. This pH range varies between indicators, but as a rule of thumb, it falls between the pKa or pKb value plus or minus one. This assumes that solutions retain their color as long as at least 10% of the other species persists. For example, if the concentration of the conjugate base is 10 times greater than the concentration of the acid, their ratio is 10:1, and consequently the pH is pKa + 1 or pKb + 1. Conversely, if a 10-fold excess of the acid occurs with respect to the base, the ratio is 1:10 and the pH is pKa − 1 or pKb − 1. For optimal accuracy, the color difference between the two species should be as clear as possible, and the narrower the pH range of the color change the better. In some indicators, such as phenolphthalein, one of the species is colorless, whereas in other indicators, such as methyl red, both species confer a color. While pH indicators work efficiently at their designated pH range, they are usually destroyed at the extreme ends of the pH scale due to undesired side reactions. Application pH indicators are frequently employed in titrations in analytical chemistry and biology to determine the extent of a chemical reaction. Because of the subjective choice (determination) of color, pH indicators are susceptible to imprecise readings. For applications requiring precise measurement of pH, a pH meter is frequently used. Sometimes, a blend of different indicators is used to achieve several smooth color changes over a wide range of pH values. These commercial indicators (e.g., universal indicator and Hydrion papers) are used when only rough knowledge of pH is necessary. For a titration, the difference between the true endpoint and the indicated endpoint is called the indicator error. Tabulated below are several common laboratory pH indicators. Indicators usually exhibit intermediate colors at pH values inside the listed transition range. For example, phenol red exhibits an orange color between pH 6.8 and pH 8.4. The transition range may shift slightly depending on the concentration of the indicator in the solution and on the temperature at which it is used. The figure on the right shows indicators with their operation range and color changes. Universal Indicator Precise pH measurement An indicator may be used to obtain quite precise measurements of pH by measuring absorbance quantitatively at two or more wavelengths. The principle can be illustrated by taking the indicator to be a simple acid, HA, which dissociates into H+ and A−. HA H+ + A− The value of the acid dissociation constant, pKa, must be known. The molar absorbances, εHA and εA− of the two species HA and A− at wavelengths λx and λy must also have been determined by previous experiment. Assuming Beer's law to be obeyed, the measured absorbances Ax and Ay at the two wavelengths are simply the sum of the absorbances due to each species. These are two equations in the two concentrations [HA] and [A−]. Once solved, the pH is obtained as If measurements are made at more than two wavelengths, the concentrations [HA] and [A−] can be calculated by linear least squares. In fact, a whole spectrum may be used for this purpose. The process is illustrated for the indicator bromocresol green. The observed spectrum (green) is the sum of the spectra of HA (gold) and of A− (blue), weighted for the concentration of the two species. When a single indicator is used, this method is limited to measurements in the pH range pKa ± 1, but this range can be extended by using mixtures of two or more indicators. Because indicators have intense absorption spectra, the indicator concentration is relatively low, and the indicator itself is assumed to have a negligible effect on pH. Equivalence point In acid-base titrations, an unfitting pH indicator may induce a color change in the indicator-containing solution before or after the actual equivalence point. As a result, different equivalence points for a solution can be concluded based on the pH indicator used. This is because the slightest color change of the indicator-containing solution suggests the equivalence point has been reached. Therefore, the most suitable pH indicator has an effective pH range, where the change in color is apparent, that encompasses the pH of the equivalence point of the solution being titrated. Naturally occurring pH indicators Many plants or plant parts contain chemicals from the naturally colored anthocyanin family of compounds. They are red in acidic solutions and blue in basic. Anthocyanins can be extracted with water or other solvents from a multitude of colored plants and plant parts, including from leaves (red cabbage); flowers (geranium, poppy, or rose petals); berries (blueberries, blackcurrant); and stems (rhubarb). Extracting anthocyanins from household plants, especially red cabbage, to form a crude pH indicator is a popular introductory chemistry demonstration. Litmus, used by alchemists in the Middle Ages and still readily available, is a naturally occurring pH indicator made from a mixture of lichen species, particularly Roccella tinctoria. The word litmus is literally from 'colored moss' in Old Norse (see Litr). The color changes between red in acid solutions and blue in alkalis. The term 'litmus test' has become a widely used metaphor for any test that purports to distinguish authoritatively between alternatives. Hydrangea macrophylla flowers can change color depending on soil acidity. In acid soils, chemical reactions occur in the soil that make aluminium available to these plants, turning the flowers blue. In alkaline soils, these reactions cannot occur and therefore aluminium is not taken up by the plant. As a result, the flowers remain pink. Another natural pH indicator is the spice turmeric. It turns yellow when exposed to acids and reddish brown when in presence of an alkalis.	Physical sciences	Chemical methods	Chemistry
85425	https://en.wikipedia.org/wiki/Metalloid	Metalloid	A metalloid is a chemical element which has a preponderance of properties in between, or that are a mixture of, those of metals and nonmetals. The word metalloid comes from the Latin metallum ("metal") and the Greek oeides ("resembling in form or appearance"). There is no standard definition of a metalloid and no complete agreement on which elements are metalloids. Despite the lack of specificity, the term remains in use in the literature. The six commonly recognised metalloids are boron, silicon, germanium, arsenic, antimony and tellurium. Five elements are less frequently so classified: carbon, aluminium, selenium, polonium and astatine. On a standard periodic table, all eleven elements are in a diagonal region of the p-block extending from boron at the upper left to astatine at lower right. Some periodic tables include a dividing line between metals and nonmetals, and the metalloids may be found close to this line. Typical metalloids have a metallic appearance, may be brittle and are only fair conductors of electricity. They can form alloys with metals, and many of their other physical properties and chemical properties are intermediate between those of metallic and nonmetallic elements. They and their compounds are used in alloys, biological agents, catalysts, flame retardants, glasses, optical storage and optoelectronics, pyrotechnics, semiconductors, and electronics. The term metalloid originally referred to nonmetals. Its more recent meaning, as a category of elements with intermediate or hybrid properties, became widespread in 1940–1960. Metalloids are sometimes called semimetals, a practice that has been discouraged, as the term semimetal has a more common usage as a specific kind of electronic band structure of a substance. In this context, only arsenic and antimony are semimetals, and commonly recognised as metalloids. Definitions Judgment-based A metalloid is an element that possesses a preponderance of properties in between, or that are a mixture of, those of metals and nonmetals, and which is therefore hard to classify as either a metal or a nonmetal. This is a generic definition that draws on metalloid attributes consistently cited in the literature. Difficulty of categorisation is a key attribute. Most elements have a mixture of metallic and nonmetallic properties, and can be classified according to which set of properties is more pronounced. Only the elements at or near the margins, lacking a sufficiently clear preponderance of either metallic or nonmetallic properties, are classified as metalloids. Boron, silicon, germanium, arsenic, antimony, and tellurium are commonly recognised as metalloids. Depending on the author, one or more from selenium, polonium, or astatine are sometimes added to the list. Boron sometimes is excluded, by itself, or with silicon. Sometimes tellurium is not regarded as a metalloid. The inclusion of antimony, polonium, and astatine as metalloids has been questioned. Other elements are occasionally classified as metalloids. These elements include hydrogen, beryllium, nitrogen, phosphorus, sulfur, zinc, gallium, tin, iodine, lead, bismuth, and radon. The term metalloid has also been used for elements that exhibit metallic lustre and electrical conductivity, and that are amphoteric, such as arsenic, antimony, vanadium, chromium, molybdenum, tungsten, tin, lead, and aluminium. The p-block metals, and nonmetals (such as carbon or nitrogen) that can form alloys with metals or modify their properties have also occasionally been considered as metalloids. Criteria-based No widely accepted definition of a metalloid exists, nor any division of the periodic table into metals, metalloids, and nonmetals; Hawkes questioned the feasibility of establishing a specific definition, noting that anomalies can be found in several attempted constructs. Classifying an element as a metalloid has been described by Sharp as "arbitrary". The number and identities of metalloids depend on what classification criteria are used. Emsley recognised four metalloids (germanium, arsenic, antimony, and tellurium); James et al. listed twelve (Emsley's plus boron, carbon, silicon, selenium, bismuth, polonium, moscovium, and livermorium). On average, seven elements are included in such lists; individual classification arrangements tend to share common ground and vary in the ill-defined margins. A single quantitative criterion such as electronegativity is commonly used, metalloids having electronegativity values from 1.8 or 1.9 to 2.2. Further examples include packing efficiency (the fraction of volume in a crystal structure occupied by atoms) and the Goldhammer–Herzfeld criterion ratio. The commonly recognised metalloids have packing efficiencies of between 34% and 41%. The Goldhammer–Herzfeld ratio, roughly equal to the cube of the atomic radius divided by the molar volume, is a simple measure of how metallic an element is, the recognised metalloids having ratios from around 0.85 to 1.1 and averaging 1.0. Other authors have relied on, for example, atomic conductance or bulk coordination number. Jones, writing on the role of classification in science, observed that "[classes] are usually defined by more than two attributes". Masterton and Slowinski used three criteria to describe the six elements commonly recognised as metalloids: metalloids have ionization energies around 200 kcal/mol (837 kJ/mol) and electronegativity values close to 2.0. They also said that metalloids are typically semiconductors, though antimony and arsenic (semimetals from a physics perspective) have electrical conductivities approaching those of metals. Selenium and polonium are suspected as not in this scheme, while astatine's status is uncertain. In this context, Vernon proposed that a metalloid is a chemical element that, in its standard state, has (a) the electronic band structure of a semiconductor or a semimetal; and (b) an intermediate first ionization potential "(say 750−1,000 kJ/mol)"; and (c) an intermediate electronegativity (1.9–2.2). Periodic table territory Location Metalloids lie on either side of the dividing line between metals and nonmetals. This can be found, in varying configurations, on some periodic tables. Elements to the lower left of the line generally display increasing metallic behaviour; elements to the upper right display increasing nonmetallic behaviour. When presented as a regular stairstep, elements with the highest critical temperature for their groups (Li, Be, Al, Ge, Sb, Po) lie just below the line. The diagonal positioning of the metalloids represents an exception to the observation that elements with similar properties tend to occur in vertical groups. A related effect can be seen in other diagonal similarities between some elements and their lower right neighbours, specifically lithium-magnesium, beryllium-aluminium, and boron-silicon. Rayner-Canham has argued that these similarities extend to carbon-phosphorus, nitrogen-sulfur, and into three d-block series. This exception arises due to competing horizontal and vertical trends in the nuclear charge. Going along a period, the nuclear charge increases with atomic number as do the number of electrons. The additional pull on outer electrons as nuclear charge increases generally outweighs the screening effect of having more electrons. With some irregularities, atoms therefore become smaller, ionization energy increases, and there is a gradual change in character, across a period, from strongly metallic, to weakly metallic, to weakly nonmetallic, to strongly nonmetallic elements. Going down a main group, the effect of increasing nuclear charge is generally outweighed by the effect of additional electrons being further away from the nucleus. Atoms generally become larger, ionization energy falls, and metallic character increases. The net effect is that the location of the metal–nonmetal transition zone shifts to the right in going down a group, and analogous diagonal similarities are seen elsewhere in the periodic table, as noted. Alternative treatments Elements bordering the metal–nonmetal dividing line are not always classified as metalloids, noting a binary classification can facilitate the establishment of rules for determining bond types between metals and nonmetals. In such cases, the authors concerned focus on one or more attributes of interest to make their classification decisions, rather than being concerned about the marginal nature of the elements in question. Their considerations may or not be made explicit and may, at times, seem arbitrary. Metalloids may be grouped with metals; or regarded as nonmetals; or treated as a sub-category of nonmetals. Other authors have suggested classifying some elements as metalloids "emphasizes that properties change gradually rather than abruptly as one moves across or down the periodic table". Some periodic tables distinguish elements that are metalloids and display no formal dividing line between metals and nonmetals. Metalloids are instead shown as occurring in a diagonal band or diffuse region. The key consideration is to explain the context for the taxonomy in use. Properties Metalloids usually look like metals but behave largely like nonmetals. Physically, they are shiny, brittle solids with intermediate to relatively good electrical conductivity and the electronic band structure of a semimetal or semiconductor. Chemically, they mostly behave as (weak) nonmetals, have intermediate ionization energies and electronegativity values, and amphoteric or weakly acidic oxides. Most of their other physical and chemical properties are intermediate in nature. Compared to metals and nonmetals Characteristic properties of metals, metalloids, and nonmetals are summarized in the table. Physical properties are listed in order of ease of determination; chemical properties run from general to specific, and then to descriptive. The above table reflects the hybrid nature of metalloids. The properties of form, appearance, and behaviour when mixed with metals are more like metals. Elasticity and general chemical behaviour are more like nonmetals. Electrical conductivity, band structure, ionization energy, electronegativity, and oxides are intermediate between the two. Common applications The focus of this section is on the recognised metalloids. Elements less often recognised as metalloids are ordinarily classified as either metals or nonmetals; some of these are included here for comparative purposes. Metalloids are too brittle to have any structural uses in their pure forms. They and their compounds are used in alloys, biological agents (toxicological, nutritional, and medicinal), catalysts, flame retardants, glasses (oxide and metallic), optical storage media and optoelectronics, pyrotechnics, semiconductors, and electronics. Alloys Writing early in the history of intermetallic compounds, the British metallurgist Cecil Desch observed that "certain non-metallic elements are capable of forming compounds of distinctly metallic character with metals, and these elements may therefore enter into the composition of alloys". He associated silicon, arsenic, and tellurium, in particular, with the alloy-forming elements. Phillips and Williams suggested that compounds of silicon, germanium, arsenic, and antimony with B metals, "are probably best classed as alloys". Among the lighter metalloids, alloys with transition metals are well-represented. Boron can form intermetallic compounds and alloys with such metals of the composition MnB, if n > 2. Ferroboron (15% boron) is used to introduce boron into steel; nickel-boron alloys are ingredients in welding alloys and case hardening compositions for the engineering industry. Alloys of silicon with iron and with aluminium are widely used by the steel and automotive industries, respectively. Germanium forms many alloys, most importantly with the coinage metals. The heavier metalloids continue the theme. Arsenic can form alloys with metals, including platinum and copper; it is also added to copper and its alloys to improve corrosion resistance and appears to confer the same benefit when added to magnesium. Antimony is well known as an alloy-former, including with the coinage metals. Its alloys include pewter (a tin alloy with up to 20% antimony) and type metal (a lead alloy with up to 25% antimony). Tellurium readily alloys with iron, as ferrotellurium (50–58% tellurium), and with copper, in the form of copper tellurium (40–50% tellurium). Ferrotellurium is used as a stabilizer for carbon in steel casting. Of the non-metallic elements less often recognised as metalloids, selenium – in the form of ferroselenium (50–58% selenium) – is used to improve the machinability of stainless steels. Biological agents All six of the elements commonly recognised as metalloids have toxic, dietary or medicinal properties. Arsenic and antimony compounds are especially toxic; boron, silicon, and possibly arsenic, are essential trace elements. Boron, silicon, arsenic, and antimony have medical applications, and germanium and tellurium are thought to have potential. Boron is used in insecticides and herbicides. It is an essential trace element. As boric acid, it has antiseptic, antifungal, and antiviral properties. Silicon is present in silatrane, a highly toxic rodenticide. Long-term inhalation of silica dust causes silicosis, a fatal disease of the lungs. Silicon is an essential trace element. Silicone gel can be applied to badly burned patients to reduce scarring. Salts of germanium are potentially harmful to humans and animals if ingested on a prolonged basis. There is interest in the pharmacological actions of germanium compounds but no licensed medicine as yet. Arsenic is notoriously poisonous and may also be an essential element in ultratrace amounts. During World War I, both sides used "arsenic-based sneezing and vomiting agents…to force enemy soldiers to remove their gas masks before firing mustard or phosgene at them in a second salvo." It has been used as a pharmaceutical agent since antiquity, including for the treatment of syphilis before the development of antibiotics. Arsenic is also a component of melarsoprol, a medicinal drug used in the treatment of human African trypanosomiasis or sleeping sickness. In 2003, arsenic trioxide (under the trade name Trisenox) was re-introduced for the treatment of acute promyelocytic leukaemia, a cancer of the blood and bone marrow. Arsenic in drinking water, which causes lung and bladder cancer, has been associated with a reduction in breast cancer mortality rates. Metallic antimony is relatively non-toxic, but most antimony compounds are poisonous. Two antimony compounds, sodium stibogluconate and stibophen, are used as antiparasitical drugs. Elemental tellurium is not considered particularly toxic; two grams of sodium tellurate, if administered, can be lethal. People exposed to small amounts of airborne tellurium exude a foul and persistent garlic-like odour. Tellurium dioxide has been used to treat seborrhoeic dermatitis; other tellurium compounds were used as antimicrobial agents before the development of antibiotics. In the future, such compounds may need to be substituted for antibiotics that have become ineffective due to bacterial resistance. Of the elements less often recognised as metalloids, beryllium and lead are noted for their toxicity; lead arsenate has been extensively used as an insecticide. Sulfur is one of the oldest of the fungicides and pesticides. Phosphorus, sulfur, zinc, selenium, and iodine are essential nutrients, and aluminium, tin, and lead may be. Sulfur, gallium, selenium, iodine, and bismuth have medicinal applications. Sulfur is a constituent of sulfonamide drugs, still widely used for conditions such as acne and urinary tract infections. Gallium nitrate is used to treat the side effects of cancer; gallium citrate, a radiopharmaceutical, facilitates imaging of inflamed body areas. Selenium sulfide is used in medicinal shampoos and to treat skin infections such as tinea versicolor. Iodine is used as a disinfectant in various forms. Bismuth is an ingredient in some antibacterials. Catalysts Boron trifluoride and trichloride are used as homogeneous catalysts in organic synthesis and electronics; the tribromide is used in the manufacture of diborane. Non-toxic boron ligands could replace toxic phosphorus ligands in some transition metal catalysts. Silica sulfuric acid (SiO2OSO3H) is used in organic reactions. Germanium dioxide is sometimes used as a catalyst in the production of PET plastic for containers; cheaper antimony compounds, such as the trioxide or triacetate, are more commonly employed for the same purpose despite concerns about antimony contamination of food and drinks. Arsenic trioxide has been used in the production of natural gas, to boost the removal of carbon dioxide, as have selenous acid and tellurous acid. Selenium acts as a catalyst in some microorganisms. Tellurium, its dioxide, and its tetrachloride are strong catalysts for air oxidation of carbon above 500 °C. Graphite oxide can be used as a catalyst in the synthesis of imines and their derivatives. Activated carbon and alumina have been used as catalysts for the removal of sulfur contaminants from natural gas. Titanium doped aluminium has been suggested as a substitute for noble metal catalysts used in the production of industrial chemicals. Flame retardants Compounds of boron, silicon, arsenic, and antimony have been used as flame retardants. Boron, in the form of borax, has been used as a textile flame retardant since at least the 18th century. Silicon compounds such as silicones, silanes, silsesquioxane, silica, and silicates, some of which were developed as alternatives to more toxic halogenated products, can considerably improve the flame retardancy of plastic materials. Arsenic compounds such as sodium arsenite or sodium arsenate are effective flame retardants for wood but have been less frequently used due to their toxicity. Antimony trioxide is a flame retardant. Aluminium hydroxide has been used as a wood-fibre, rubber, plastic, and textile flame retardant since the 1890s. Apart from aluminium hydroxide, use of phosphorus based flame-retardants – in the form of, for example, organophosphates – now exceeds that of any of the other main retardant types. These employ boron, antimony, or halogenated hydrocarbon compounds. Glass formation The oxides B2O3, SiO2, GeO2, As2O3, and Sb2O3 readily form glasses. TeO2 forms a glass but this requires a "heroic quench rate" or the addition of an impurity; otherwise the crystalline form results. These compounds are used in chemical, domestic, and industrial glassware and optics. Boron trioxide is used as a glass fibre additive, and is also a component of borosilicate glass, widely used for laboratory glassware and domestic ovenware for its low thermal expansion. Most ordinary glassware is made from silicon dioxide. Germanium dioxide is used as a glass fibre additive, as well as in infrared optical systems. Arsenic trioxide is used in the glass industry as a decolourizing and fining agent (for the removal of bubbles), as is antimony trioxide. Tellurium dioxide finds application in laser and nonlinear optics. Amorphous metallic glasses are generally most easily prepared if one of the components is a metalloid or "near metalloid" such as boron, carbon, silicon, phosphorus or germanium. Aside from thin films deposited at very low temperatures, the first known metallic glass was an alloy of composition Au75Si25 reported in 1960. A metallic glass having a strength and toughness not previously seen, of composition Pd82.5P6Si9.5Ge2, was reported in 2011. Phosphorus, selenium, and lead, which are less often recognised as metalloids, are also used in glasses. Phosphate glass has a substrate of phosphorus pentoxide (P2O5), rather than the silica (SiO2) of conventional silicate glasses. It is used, for example, to make sodium lamps. Selenium compounds can be used both as decolourising agents and to add a red colour to glass. Decorative glassware made of traditional lead glass contains at least 30% lead(II) oxide (PbO); lead glass used for radiation shielding may have up to 65% PbO. Lead-based glasses have also been extensively used in electronic components, enamelling, sealing and glazing materials, and solar cells. Bismuth based oxide glasses have emerged as a less toxic replacement for lead in many of these applications. Optical storage and optoelectronics Varying compositions of GeSbTe ("GST alloys") and Ag- and In- doped Sb2Te ("AIST alloys"), being examples of phase-change materials, are widely used in rewritable optical discs and phase-change memory devices. By applying heat, they can be switched between amorphous (glassy) and crystalline states. The change in optical and electrical properties can be used for information storage purposes. Future applications for GeSbTe may include, "ultrafast, entirely solid-state displays with nanometre-scale pixels, semi-transparent 'smart' glasses, 'smart' contact lenses, and artificial retina devices." Pyrotechnics The recognised metalloids have either pyrotechnic applications or associated properties. Boron and silicon are commonly encountered; they act somewhat like metal fuels. Boron is used in pyrotechnic initiator compositions (for igniting other hard-to-start compositions), and in delay compositions that burn at a constant rate. Boron carbide has been identified as a possible replacement for more toxic barium or hexachloroethane mixtures in smoke munitions, signal flares, and fireworks. Silicon, like boron, is a component of initiator and delay mixtures. Doped germanium can act as a variable speed thermite fuel. Arsenic trisulfide As2S3 was used in old naval signal lights; in fireworks to make white stars; in yellow smoke screen mixtures; and in initiator compositions. Antimony trisulfide Sb2S3 is found in white-light fireworks and in flash and sound mixtures. Tellurium has been used in delay mixtures and in blasting cap initiator compositions. Carbon, aluminium, phosphorus, and selenium continue the theme. Carbon, in black powder, is a constituent of fireworks rocket propellants, bursting charges, and effects mixtures, and military delay fuses and igniters. Aluminium is a common pyrotechnic ingredient, and is widely employed for its capacity to generate light and heat, including in thermite mixtures. Phosphorus can be found in smoke and incendiary munitions, paper caps used in toy guns, and party poppers. Selenium has been used in the same way as tellurium. Semiconductors and electronics All the elements commonly recognised as metalloids (or their compounds) have been used in the semiconductor or solid-state electronic industries. Some properties of boron have limited its use as a semiconductor. It has a high melting point, single crystals are relatively hard to obtain, and introducing and retaining controlled impurities is difficult. Silicon is the leading commercial semiconductor; it forms the basis of modern electronics (including standard solar cells) and information and communication technologies. This was despite the study of semiconductors, early in the 20th century, having been regarded as the "physics of dirt" and not deserving of close attention. Germanium has largely been replaced by silicon in semiconducting devices, being cheaper, more resilient at higher operating temperatures, and easier to work during the microelectronic fabrication process. Germanium is still a constituent of semiconducting silicon-germanium "alloys" and these have been growing in use, particularly for wireless communication devices; such alloys exploit the higher carrier mobility of germanium. The synthesis of gram-scale quantities of semiconducting germanane was reported in 2013. This consists of one-atom thick sheets of hydrogen-terminated germanium atoms, analogous to graphane. It conducts electrons more than ten times faster than silicon and five times faster than germanium, and is thought to have potential for optoelectronic and sensing applications. The development of a germanium-wire based anode that more than doubles the capacity of lithium-ion batteries was reported in 2014. In the same year, Lee et al. reported that defect-free crystals of graphene large enough to have electronic uses could be grown on, and removed from, a germanium substrate. Arsenic and antimony are not semiconductors in their standard states. Both form type III-V semiconductors (such as GaAs, AlSb or GaInAsSb) in which the average number of valence electrons per atom is the same as that of Group 14 elements, but they have direct band gaps. These compounds are preferred for optical applications. Antimony nanocrystals may enable lithium-ion batteries to be replaced by more powerful sodium ion batteries. Tellurium, which is a semiconductor in its standard state, is used mainly as a component in type II/VI semiconducting-chalcogenides; these have applications in electro-optics and electronics. Cadmium telluride (CdTe) is used in solar modules for its high conversion efficiency, low manufacturing costs, and large band gap of 1.44 eV, letting it absorb a wide range of wavelengths. Bismuth telluride (Bi2Te3), alloyed with selenium and antimony, is a component of thermoelectric devices used for refrigeration or portable power generation. Five metalloids – boron, silicon, germanium, arsenic, and antimony – can be found in cell phones (along with at least 39 other metals and nonmetals). Tellurium is expected to find such use. Of the less often recognised metalloids, phosphorus, gallium (in particular) and selenium have semiconductor applications. Phosphorus is used in trace amounts as a dopant for n-type semiconductors. The commercial use of gallium compounds is dominated by semiconductor applications – in integrated circuits, cell phones, laser diodes, light-emitting diodes, photodetectors, and solar cells. Selenium is used in the production of solar cells and in high-energy surge protectors. Boron, silicon, germanium, antimony, and tellurium, as well as heavier metals and metalloids such as Sm, Hg, Tl, Pb, Bi, and Se, can be found in topological insulators. These are alloys or compounds which, at ultracold temperatures or room temperature (depending on their composition), are metallic conductors on their surfaces but insulators through their interiors. Cadmium arsenide Cd3As2, at about 1 K, is a Dirac-semimetal – a bulk electronic analogue of graphene – in which electrons travel effectively as massless particles. These two classes of material are thought to have potential quantum computing applications. Nomenclature and history Derivation and other names Several names are sometimes used synonymously although some of these have other meanings that are not necessarily interchangeable: amphoteric element, boundary element, half-way element, near metal, meta-metal, semiconductor, semimetal and submetal. "Amphoteric element" is sometimes used more broadly to include transition metals capable of forming oxyanions, such as chromium and manganese. "Meta-metal" is sometimes used instead to refer to certain metals (Be, Zn, Cd, Hg, In, Tl, β-Sn, Pb) located just to the left of the metalloids on standard periodic tables. These metals tend to have distorted crystalline structures, electrical conductivity values at the lower end of those of metals, and amphoteric (weakly basic) oxides. The names amphoteric element and semiconductor are problematic as some elements referred to as metalloids do not show marked amphoteric behaviour (bismuth, for example) or semiconductivity (polonium) in their most stable forms. Origin and usage The origin and usage of the term metalloid is convoluted. The "Manual of Metalloids" published in 1864 divided all elements into either metals or metalloids. Earlier usage in mineralogy, to describe a mineral having a metallic appearance, can be sourced to as early as 1800. Since the mid-20th century it has been used to refer to intermediate or borderline chemical elements. The International Union of Pure and Applied Chemistry (IUPAC) previously recommended abandoning the term metalloid, and suggested using the term semimetal instead. Use of this latter term has more recently been discouraged by Atkins et al. as it has a more common meaning that refers to the electronic band structure of a substance rather than the overall classification of an element. The most recent IUPAC publications on nomenclature and terminology do not include any recommendations on the usage of the terms metalloid or semimetal. Elements commonly recognised as metalloids Properties noted in this section refer to the elements in their most thermodynamically stable forms under ambient conditions. Boron Pure boron is a shiny, silver-grey crystalline solid. It is less dense than aluminium (2.34 vs. 2.70 g/cm3), and is hard and brittle. It is barely reactive under normal conditions, except for attack by fluorine, and has a melting point of 2076 °C (cf. steel ~1370 °C). Boron is a semiconductor; its room temperature electrical conductivity is 1.5 × 10−6 S•cm−1 (about 200 times less than that of tap water) and it has a band gap of about 1.56 eV. Mendeleev commented that, "Boron appears in a free state in several forms which are intermediate between the metals and the nonmmetals." The structural chemistry of boron is dominated by its small atomic size, and relatively high ionization energy. With only three valence electrons per boron atom, simple covalent bonding cannot fulfil the octet rule. Metallic bonding is the usual result among the heavier congenors of boron but this generally requires low ionization energies. Instead, because of its small size and high ionization energies, the basic structural unit of boron (and nearly all of its allotropes) is the icosahedral B12 cluster. Of the 36 electrons associated with 12 boron atoms, 26 reside in 13 delocalized molecular orbitals; the other 10 electrons are used to form two- and three-centre covalent bonds between icosahedra. The same motif can be seen, as are deltahedral variants or fragments, in metal borides and hydride derivatives, and in some halides. The bonding in boron has been described as being characteristic of behaviour intermediate between metals and nonmetallic covalent network solids (such as diamond). The energy required to transform B, C, N, Si, and P from nonmetallic to metallic states has been estimated as 30, 100, 240, 33, and 50 kJ/mol, respectively. This indicates the proximity of boron to the metal-nonmetal borderline. Most of the chemistry of boron is nonmetallic in nature. Unlike its heavier congeners, it is not known to form a simple B3+ or hydrated [B(H2O)4]3+ cation. The small size of the boron atom enables the preparation of many interstitial alloy-type borides. Analogies between boron and transition metals have been noted in the formation of complexes, and adducts (for example, BH3 + CO →BH3CO and, similarly, Fe(CO)4 + CO →Fe(CO)5), as well as in the geometric and electronic structures of cluster species such as [B6H6]2− and [Ru6(CO)18]2−. The aqueous chemistry of boron is characterised by the formation of many different polyborate anions. Given its high charge-to-size ratio, boron bonds covalently in nearly all of its compounds; the exceptions are the borides as these include, depending on their composition, covalent, ionic, and metallic bonding components. Simple binary compounds, such as boron trichloride are Lewis acids as the formation of three covalent bonds leaves a hole in the octet which can be filled by an electron-pair donated by a Lewis base. Boron has a strong affinity for oxygen and a duly extensive borate chemistry. The oxide B2O3 is polymeric in structure, weakly acidic, and a glass former. Organometallic compounds of boron have been known since the 19th century (see organoboron chemistry). Silicon Silicon is a crystalline solid with a blue-grey metallic lustre. Like boron, it is less dense (at 2.33 g/cm3) than aluminium, and is hard and brittle. It is a relatively unreactive element. According to Rochow, the massive crystalline form (especially if pure) is "remarkably inert to all acids, including hydrofluoric". Less pure silicon, and the powdered form, are variously susceptible to attack by strong or heated acids, as well as by steam and fluorine. Silicon dissolves in hot aqueous alkalis with the evolution of hydrogen, as do metals such as beryllium, aluminium, zinc, gallium or indium. It melts at 1414 °C. Silicon is a semiconductor with an electrical conductivity of 10−4 S•cm−1 and a band gap of about 1.11 eV. When it melts, silicon becomes a reasonable metal with an electrical conductivity of 1.0–1.3 × 104 S•cm−1, similar to that of liquid mercury. The chemistry of silicon is generally nonmetallic (covalent) in nature. It is not known to form a cation. Silicon can form alloys with metals such as iron and copper. It shows fewer tendencies to anionic behaviour than ordinary nonmetals. Its solution chemistry is characterised by the formation of oxyanions. The high strength of the silicon–oxygen bond dominates the chemical behaviour of silicon. Polymeric silicates, built up by tetrahedral SiO4 units sharing their oxygen atoms, are the most abundant and important compounds of silicon. The polymeric borates, comprising linked trigonal and tetrahedral BO3 or BO4 units, are built on similar structural principles. The oxide SiO2 is polymeric in structure, weakly acidic, and a glass former. Traditional organometallic chemistry includes the carbon compounds of silicon (see organosilicon). Germanium Germanium is a shiny grey-white solid. It has a density of 5.323 g/cm3 and is hard and brittle. It is mostly unreactive at room temperature but is slowly attacked by hot concentrated sulfuric or nitric acid. Germanium also reacts with molten caustic soda to yield sodium germanate Na2GeO3 and hydrogen gas. It melts at 938 °C. Germanium is a semiconductor with an electrical conductivity of around 2 × 10−2 S•cm−1 and a band gap of 0.67 eV. Liquid germanium is a metallic conductor, with an electrical conductivity similar to that of liquid mercury. Most of the chemistry of germanium is characteristic of a nonmetal. Whether or not germanium forms a cation is unclear, aside from the reported existence of the Ge2+ ion in a few esoteric compounds. It can form alloys with metals such as aluminium and gold. It shows fewer tendencies to anionic behaviour than ordinary nonmetals. Its solution chemistry is characterised by the formation of oxyanions. Germanium generally forms tetravalent (IV) compounds, and it can also form less stable divalent (II) compounds, in which it behaves more like a metal. Germanium analogues of all of the major types of silicates have been prepared. The metallic character of germanium is also suggested by the formation of various oxoacid salts. A phosphate [(HPO4)2Ge·H2O] and highly stable trifluoroacetate Ge(OCOCF3)4 have been described, as have Ge2(SO4)2, Ge(ClO4)4 and GeH2(C2O4)3. The oxide GeO2 is polymeric, amphoteric, and a glass former. The dioxide is soluble in acidic solutions (the monoxide GeO, is even more so), and this is sometimes used to classify germanium as a metal. Up to the 1930s germanium was considered to be a poorly conducting metal; it has occasionally been classified as a metal by later writers. As with all the elements commonly recognised as metalloids, germanium has an established organometallic chemistry (see Organogermanium chemistry). Arsenic Arsenic is a grey, metallic looking solid. It has a density of 5.727 g/cm3 and is brittle, and moderately hard (more than aluminium; less than iron). It is stable in dry air but develops a golden bronze patina in moist air, which blackens on further exposure. Arsenic is attacked by nitric acid and concentrated sulfuric acid. It reacts with fused caustic soda to give the arsenate Na3AsO3 and hydrogen gas. Arsenic sublimes at 615 °C. The vapour is lemon-yellow and smells like garlic. Arsenic only melts under a pressure of 38.6 atm, at 817 °C. It is a semimetal with an electrical conductivity of around 3.9 × 104 S•cm−1 and a band overlap of 0.5 eV. Liquid arsenic is a semiconductor with a band gap of 0.15 eV. The chemistry of arsenic is predominately nonmetallic. Whether or not arsenic forms a cation is unclear. Its many metal alloys are mostly brittle. It shows fewer tendencies to anionic behaviour than ordinary nonmetals. Its solution chemistry is characterised by the formation of oxyanions. Arsenic generally forms compounds in which it has an oxidation state of +3 or +5. The halides, and the oxides and their derivatives are illustrative examples. In the trivalent state, arsenic shows some incipient metallic properties. The halides are hydrolysed by water but these reactions, particularly those of the chloride, are reversible with the addition of a hydrohalic acid. The oxide is acidic but, as noted below, (weakly) amphoteric. The higher, less stable, pentavalent state has strongly acidic (nonmetallic) properties. Compared to phosphorus, the stronger metallic character of arsenic is indicated by the formation of oxoacid salts such as AsPO4, As2(SO4)3 and arsenic acetate As(CH3COO)3. The oxide As2O3 is polymeric, amphoteric, and a glass former. Arsenic has an extensive organometallic chemistry (see Organoarsenic chemistry). Antimony Antimony is a silver-white solid with a blue tint and a brilliant lustre. It has a density of 6.697 g/cm3 and is brittle, and moderately hard (more so than arsenic; less so than iron; about the same as copper). It is stable in air and moisture at room temperature. It is attacked by concentrated nitric acid, yielding the hydrated pentoxide Sb2O5. Aqua regia gives the pentachloride SbCl5 and hot concentrated sulfuric acid results in the sulfate Sb2(SO4)3. It is not affected by molten alkali. Antimony is capable of displacing hydrogen from water, when heated: 2 Sb + 3 H2O → Sb2O3 + 3 H2. It melts at 631 °C. Antimony is a semimetal with an electrical conductivity of around 3.1 × 104 S•cm−1 and a band overlap of 0.16 eV. Liquid antimony is a metallic conductor with an electrical conductivity of around 5.3 × 104 S•cm−1. Most of the chemistry of antimony is characteristic of a nonmetal. Antimony has some definite cationic chemistry, SbO+ and Sb(OH)2+ being present in acidic aqueous solution; the compound Sb8(GaCl4)2, which contains the homopolycation, Sb82+, was prepared in 2004. It can form alloys with one or more metals such as aluminium, iron, nickel, copper, zinc, tin, lead, and bismuth. Antimony has fewer tendencies to anionic behaviour than ordinary nonmetals. Its solution chemistry is characterised by the formation of oxyanions. Like arsenic, antimony generally forms compounds in which it has an oxidation state of +3 or +5. The halides, and the oxides and their derivatives are illustrative examples. The +5 state is less stable than the +3, but relatively easier to attain than with arsenic. This is explained by the poor shielding afforded the arsenic nucleus by its 3d10 electrons. In comparison, the tendency of antimony (being a heavier atom) to oxidize more easily partially offsets the effect of its 4d10 shell. Tripositive antimony is amphoteric; pentapositive antimony is (predominately) acidic. Consistent with an increase in metallic character down group 15, antimony forms salts including an acetate Sb(CH3CO2)3, phosphate SbPO4, sulfate Sb2(SO4)3 and perchlorate Sb(ClO4)3. The otherwise acidic pentoxide Sb2O5 shows some basic (metallic) behaviour in that it can be dissolved in very acidic solutions, with the formation of the oxycation SbO. The oxide Sb2O3 is polymeric, amphoteric, and a glass former. Antimony has an extensive organometallic chemistry (see Organoantimony chemistry). Tellurium Tellurium is a silvery-white shiny solid. It has a density of 6.24 g/cm3, is brittle, and is the softest of the commonly recognised metalloids, being marginally harder than sulfur. Large pieces of tellurium are stable in air. The finely powdered form is oxidized by air in the presence of moisture. Tellurium reacts with boiling water, or when freshly precipitated even at 50 °C, to give the dioxide and hydrogen: Te + 2 H2O → TeO2 + 2 H2. It reacts (to varying degrees) with nitric, sulfuric, and hydrochloric acids to give compounds such as the sulfoxide TeSO3 or tellurous acid H2TeO3, the basic nitrate (Te2O4H)+(NO3)−, or the oxide sulfate Te2O3(SO4). It dissolves in boiling alkalis, to give the tellurite and telluride: 3 Te + 6 KOH = K2TeO3 + 2 K2Te + 3 H2O, a reaction that proceeds or is reversible with increasing or decreasing temperature. At higher temperatures tellurium is sufficiently plastic to extrude. It melts at 449.51 °C. Crystalline tellurium has a structure consisting of parallel infinite spiral chains. The bonding between adjacent atoms in a chain is covalent, but there is evidence of a weak metallic interaction between the neighbouring atoms of different chains. Tellurium is a semiconductor with an electrical conductivity of around 1.0 S•cm−1 and a band gap of 0.32 to 0.38 eV. Liquid tellurium is a semiconductor, with an electrical conductivity, on melting, of around 1.9 × 103 S•cm−1. Superheated liquid tellurium is a metallic conductor. Most of the chemistry of tellurium is characteristic of a nonmetal. It shows some cationic behaviour. The dioxide dissolves in acid to yield the trihydroxotellurium(IV) Te(OH)3+ ion; the red Te42+ and yellow-orange Te62+ ions form when tellurium is oxidized in fluorosulfuric acid (HSO3F), or liquid sulfur dioxide (SO2), respectively. It can form alloys with aluminium, silver, and tin. Tellurium shows fewer tendencies to anionic behaviour than ordinary nonmetals. Its solution chemistry is characterised by the formation of oxyanions. Tellurium generally forms compounds in which it has an oxidation state of −2, +4 or +6. The +4 state is the most stable. Tellurides of composition XxTey are easily formed with most other elements and represent the most common tellurium minerals. Nonstoichiometry is pervasive, especially with transition metals. Many tellurides can be regarded as metallic alloys. The increase in metallic character evident in tellurium, as compared to the lighter chalcogens, is further reflected in the reported formation of various other oxyacid salts, such as a basic selenate 2TeO2·SeO3 and an analogous perchlorate and periodate 2TeO2·HXO4. Tellurium forms a polymeric, amphoteric, glass-forming oxide TeO2. It is a "conditional" glass-forming oxide – it forms a glass with a very small amount of additive. Tellurium has an extensive organometallic chemistry (see Organotellurium chemistry). Elements less commonly recognised as metalloids Carbon Carbon is ordinarily classified as a nonmetal but has some metallic properties and is occasionally classified as a metalloid. Hexagonal graphitic carbon (graphite) is the most thermodynamically stable allotrope of carbon under ambient conditions. It has a lustrous appearance and is a fairly good electrical conductor. Graphite has a layered structure. Each layer consists of carbon atoms bonded to three other carbon atoms in a hexagonal lattice arrangement. The layers are stacked together and held loosely by van der Waals forces and delocalized valence electrons. Like a metal, the conductivity of graphite in the direction of its planes decreases as the temperature is raised; it has the electronic band structure of a semimetal. The allotropes of carbon, including graphite, can accept foreign atoms or compounds into their structures via substitution, intercalation, or doping. The resulting materials are sometimes referred to as "carbon alloys". Carbon can form ionic salts, including a hydrogen sulfate, perchlorate, and nitrate (CX−.2HX, where X = HSO4, ClO4; and CNO.3HNO3). In organic chemistry, carbon can form complex cationstermed carbocationsin which the positive charge is on the carbon atom; examples are and , and their derivatives. Graphite is an established solid lubricant and behaves as a semiconductor in a direction perpendicular to its planes. Most of its chemistry is nonmetallic; it has a relatively high ionization energy and, compared to most metals, a relatively high electronegativity. Carbon can form anions such as C4− (methanide), C (acetylide), and C (sesquicarbide or allylenide), in compounds with metals of main groups 1–3, and with the lanthanides and actinides. Its oxide CO2 forms carbonic acid H2CO3. Aluminium Aluminium is ordinarily classified as a metal. It is lustrous, malleable and ductile, and has high electrical and thermal conductivity. Like most metals it has a close-packed crystalline structure, and forms a cation in aqueous solution. It has some properties that are unusual for a metal; taken together, these are sometimes used as a basis to classify aluminium as a metalloid. Its crystalline structure shows some evidence of directional bonding. Aluminium bonds covalently in most compounds. The oxide Al2O3 is amphoteric and a conditional glass-former. Aluminium can form anionic aluminates, such behaviour being considered nonmetallic in character. Classifying aluminium as a metalloid has been disputed given its many metallic properties. It is therefore, arguably, an exception to the mnemonic that elements adjacent to the metal–nonmetal dividing line are metalloids. Stott labels aluminium as a weak metal. It has the physical properties of a metal but some of the chemical properties of a nonmetal. Steele notes the paradoxical chemical behaviour of aluminium: "It resembles a weak metal in its amphoteric oxide and in the covalent character of many of its compounds ... Yet it is a highly electropositive metal ... [with] a high negative electrode potential". Moody says that, "aluminium is on the 'diagonal borderland' between metals and non-metals in the chemical sense." Selenium Selenium shows borderline metalloid or nonmetal behaviour. Its most stable form, the grey trigonal allotrope, is sometimes called "metallic" selenium because its electrical conductivity is several orders of magnitude greater than that of the red monoclinic form. The metallic character of selenium is further shown by its lustre, and its crystalline structure, which is thought to include weakly "metallic" interchain bonding. Selenium can be drawn into thin threads when molten and viscous. It shows reluctance to acquire "the high positive oxidation numbers characteristic of nonmetals". It can form cyclic polycations (such as Se) when dissolved in oleums (an attribute it shares with sulfur and tellurium), and a hydrolysed cationic salt in the form of trihydroxoselenium(IV) perchlorate [Se(OH)3]+·ClO. The nonmetallic character of selenium is shown by its brittleness and the low electrical conductivity (~10−9 to 10−12 S•cm−1) of its highly purified form. This is comparable to or less than that of bromine (7.95 S•cm−1), a nonmetal. Selenium has the electronic band structure of a semiconductor and retains its semiconducting properties in liquid form. It has a relatively high electronegativity (2.55 revised Pauling scale). Its reaction chemistry is mainly that of its nonmetallic anionic forms Se2−, SeO and SeO. Selenium is commonly described as a metalloid in the environmental chemistry literature. It moves through the aquatic environment similarly to arsenic and antimony; its water-soluble salts, in higher concentrations, have a similar toxicological profile to that of arsenic. Polonium Polonium is "distinctly metallic" in some ways. Both of its allotropic forms are metallic conductors. It is soluble in acids, forming the rose-coloured Po2+ cation and displacing hydrogen: Po + 2 H+ → Po2+ + H2. Many polonium salts are known. The oxide PoO2 is predominantly basic in nature. Polonium is a reluctant oxidizing agent, unlike its lightest congener oxygen: highly reducing conditions are required for the formation of the Po2− anion in aqueous solution. Whether polonium is ductile or brittle is unclear. It is predicted to be ductile based on its calculated elastic constants. It has a simple cubic crystalline structure. Such a structure has few slip systems and "leads to very low ductility and hence low fracture resistance". Polonium shows nonmetallic character in its halides, and by the existence of polonides. The halides have properties generally characteristic of nonmetal halides (being volatile, easily hydrolyzed, and soluble in organic solvents). Many metal polonides, obtained by heating the elements together at 500–1,000 °C, and containing the Po2− anion, are also known. Astatine As a halogen, astatine tends to be classified as a nonmetal. It has some marked metallic properties and is sometimes instead classified as either a metalloid or (less often) as a metal. Immediately following its production in 1940, early investigators considered it a metal. In 1949 it was called the most noble (difficult to reduce) nonmetal as well as being a relatively noble (difficult to oxidize) metal. In 1950 astatine was described as a halogen and (therefore) a reactive nonmetal. In 2013, on the basis of relativistic modelling, astatine was predicted to be a monatomic metal, with a face-centred cubic crystalline structure. Several authors have commented on the metallic nature of some of the properties of astatine. Since iodine is a semiconductor in the direction of its planes, and since the halogens become more metallic with increasing atomic number, it has been presumed that astatine would be a metal if it could form a condensed phase. Astatine may be metallic in the liquid state on the basis that elements with an enthalpy of vaporization (∆Hvap) greater than ~42 kJ/mol are metallic when liquid. Such elements include boron, silicon, germanium, antimony, selenium, and tellurium. Estimated values for ∆Hvap of diatomic astatine are 50 kJ/mol or higher; diatomic iodine, with a ∆Hvap of 41.71, falls just short of the threshold figure. "Like typical metals, it [astatine] is precipitated by hydrogen sulfide even from strongly acid solutions and is displaced in a free form from sulfate solutions; it is deposited on the cathode on electrolysis." Further indications of a tendency for astatine to behave like a (heavy) metal are: "... the formation of pseudohalide compounds ... complexes of astatine cations ... complex anions of trivalent astatine ... as well as complexes with a variety of organic solvents". It has also been argued that astatine demonstrates cationic behaviour, by way of stable At+ and AtO+ forms, in strongly acidic aqueous solutions. Some of astatine's reported properties are nonmetallic. It has been extrapolated to have the narrow liquid range ordinarily associated with nonmetals (mp 302 °C; bp 337 °C), although experimental indications suggest a lower boiling point of about 230±3 °C. Batsanov gives a calculated band gap energy for astatine of 0.7 eV; this is consistent with nonmetals (in physics) having separated valence and conduction bands and thereby being either semiconductors or insulators. The chemistry of astatine in aqueous solution is mainly characterised by the formation of various anionic species. Most of its known compounds resemble those of iodine, which is a halogen and a nonmetal. Such compounds include astatides (XAt), astatates (XAtO3), and monovalent interhalogen compounds. Restrepo et al. reported that astatine appeared to be more polonium-like than halogen-like. They did so on the basis of detailed comparative studies of the known and interpolated properties of 72 elements. Related concepts Near metalloids In the periodic table, some of the elements adjacent to the commonly recognised metalloids, although usually classified as either metals or nonmetals, are occasionally referred to as near-metalloids or noted for their metalloidal character. To the left of the metal–nonmetal dividing line, such elements include gallium, tin and bismuth. They show unusual packing structures, marked covalent chemistry (molecular or polymeric), and amphoterism. To the right of the dividing line are carbon, phosphorus, selenium and iodine. They exhibit metallic lustre, semiconducting properties and bonding or valence bands with delocalized character. This applies to their most thermodynamically stable forms under ambient conditions: carbon as graphite; phosphorus as black phosphorus; and selenium as grey selenium. Allotropes Different crystalline forms of an element are called allotropes. Some allotropes, particularly those of elements located (in periodic table terms) alongside or near the notional dividing line between metals and nonmetals, exhibit more pronounced metallic, metalloidal or nonmetallic behaviour than others. The existence of such allotropes can complicate the classification of the elements involved. Tin, for example, has two allotropes: tetragonal "white" β-tin and cubic "grey" α-tin. White tin is a very shiny, ductile and malleable metal. It is the stable form at or above room temperature and has an electrical conductivity of 9.17 × 104 S·cm−1 (~1/6th that of copper). Grey tin usually has the appearance of a grey micro-crystalline powder, and can also be prepared in brittle semi-lustrous crystalline or polycrystalline forms. It is the stable form below 13.2 °C and has an electrical conductivity of between (2–5) × 102 S·cm−1 (~1/250th that of white tin). Grey tin has the same crystalline structure as that of diamond. It behaves as a semiconductor (as if it had a band gap of 0.08 eV), but has the electronic band structure of a semimetal. It has been referred to as either a very poor metal, a metalloid, a nonmetal or a near metalloid. The diamond allotrope of carbon is clearly nonmetallic, being translucent and having a low electrical conductivity of 10−14 to 10−16 S·cm−1. Graphite has an electrical conductivity of 3 × 104 S·cm−1, a figure more characteristic of a metal. Phosphorus, sulfur, arsenic, selenium, antimony, and bismuth also have less stable allotropes that display different behaviours. Abundance, extraction, and cost Abundance The table gives crustal abundances of the elements commonly to rarely recognised as metalloids. Some other elements are included for comparison: oxygen and xenon (the most and least abundant elements with stable isotopes); iron and the coinage metals copper, silver, and gold; and rhenium, the least abundant stable metal (aluminium is normally the most abundant metal). Various abundance estimates have been published; these often disagree to some extent. Extraction The recognised metalloids can be obtained by chemical reduction of either their oxides or their sulfides. Simpler or more complex extraction methods may be employed depending on the starting form and economic factors. Boron is routinely obtained by reducing the trioxide with magnesium: B2O3 + 3 Mg → 2 B + 3MgO; after secondary processing the resulting brown powder has a purity of up to 97%. Boron of higher purity (> 99%) is prepared by heating volatile boron compounds, such as BCl3 or BBr3, either in a hydrogen atmosphere (2 BX3 + 3 H2 → 2 B + 6 HX) or to the point of thermal decomposition. Silicon and germanium are obtained from their oxides by heating the oxide with carbon or hydrogen: SiO2 + C → Si + CO2; GeO2 + 2 H2 → Ge + 2 H2O. Arsenic is isolated from its pyrite (FeAsS) or arsenical pyrite (FeAs2) by heating; alternatively, it can be obtained from its oxide by reduction with carbon: 2 As2O3 + 3 C → 2 As + 3 CO2. Antimony is derived from its sulfide by reduction with iron: Sb2S3 → 2 Sb + 3 FeS. Tellurium is prepared from its oxide by dissolving it in aqueous NaOH, yielding tellurite, then by electrolytic reduction: TeO2 + 2 NaOH → Na2TeO3 + H2O; Na2TeO3 + H2O → Te + 2 NaOH + O2. Another option is reduction of the oxide by roasting with carbon: TeO2 + C → Te + CO2. Production methods for the elements less frequently recognised as metalloids involve natural processing, electrolytic or chemical reduction, or irradiation. Carbon (as graphite) occurs naturally and is extracted by crushing the parent rock and floating the lighter graphite to the surface. Aluminium is extracted by dissolving its oxide Al2O3 in molten cryolite Na3AlF6 and then by high temperature electrolytic reduction. Selenium is produced by roasting the coinage metal selenides X2Se (X = Cu, Ag, Au) with soda ash to give the selenite: X2Se + O2 + Na2CO3 → Na2SeO3 + 2 X + CO2; the selenide is neutralized by sulfuric acid H2SO4 to give selenous acid H2SeO3; this is reduced by bubbling with SO2 to yield elemental selenium. Polonium and astatine are produced in minute quantities by irradiating bismuth. Cost The recognised metalloids and their closer neighbours mostly cost less than silver; only polonium and astatine are more expensive than gold, on account of their significant radioactivity. As of 5 April 2014, prices for small samples (up to 100 g) of silicon, antimony and tellurium, and graphite, aluminium and selenium, average around one third the cost of silver (US$1.5 per gram or about $45 an ounce). Boron, germanium, and arsenic samples average about three-and-a-half times the cost of silver. Polonium is available for about $100 per microgram. Zalutsky and Pruszynski estimate a similar cost for producing astatine. Prices for the applicable elements traded as commodities tend to range from two to three times cheaper than the sample price (Ge), to nearly three thousand times cheaper (As).	Physical sciences	Chemical element groups	null
85631	https://en.wikipedia.org/wiki/Augmented%20reality	Augmented reality	Augmented reality (AR) is an interactive experience that combines the real world and computer-generated 3D content. The content can span multiple sensory modalities, including visual, auditory, haptic, somatosensory and olfactory. AR can be defined as a system that incorporates three basic features: a combination of real and virtual worlds, real-time interaction, and accurate 3D registration of virtual and real objects. The overlaid sensory information can be constructive (i.e. additive to the natural environment), or destructive (i.e. masking of the natural environment). As such, it is one of the key technologies in the reality-virtuality continuum. This experience is seamlessly interwoven with the physical world such that it is perceived as an immersive aspect of the real environment. In this way, augmented reality alters one's ongoing perception of a real-world environment, whereas virtual reality completely replaces the user's real-world environment with a simulated one. Augmented reality is largely synonymous with mixed reality. There is also overlap in terminology with extended reality and computer-mediated reality. The primary value of augmented reality is the manner in which components of the digital world blend into a person's perception of the real world, not as a simple display of data, but through the integration of immersive sensations, which are perceived as natural parts of an environment. The earliest functional AR systems that provided immersive mixed reality experiences for users were invented in the early 1990s, starting with the Virtual Fixtures system developed at the U.S. Air Force's Armstrong Laboratory in 1992. Commercial augmented reality experiences were first introduced in entertainment and gaming businesses. Subsequently, augmented reality applications have spanned commercial industries such as education, communications, medicine, and entertainment. In education, content may be accessed by scanning or viewing an image with a mobile device or by using markerless AR techniques. Augmented reality can be used to enhance natural environments or situations and offers perceptually enriched experiences. With the help of advanced AR technologies (e.g. adding computer vision, incorporating AR cameras into smartphone applications, and object recognition) the information about the surrounding real world of the user becomes interactive and digitally manipulated. Information about the environment and its objects is overlaid on the real world. This information can be virtual. Augmented Reality is any experience which is artificial and which adds to the already existing reality. or real, e.g. seeing other real sensed or measured information such as electromagnetic radio waves overlaid in exact alignment with where they actually are in space. Augmented reality also has a lot of potential in the gathering and sharing of tacit knowledge. Augmentation techniques are typically performed in real-time and in semantic contexts with environmental elements. Immersive perceptual information is sometimes combined with supplemental information like scores over a live video feed of a sporting event. This combines the benefits of both augmented reality technology and heads up display technology (HUD). Comparison with virtual reality In virtual reality (VR), the users' perception is completely computer-generated, whereas with augmented reality (AR), it is partially generated and partially from the real world. For example, in architecture, VR can be used to create a walk-through simulation of the inside of a new building; and AR can be used to show a building's structures and systems super-imposed on a real-life view. Another example is through the use of utility applications. Some AR applications, such as Augment, enable users to apply digital objects into real environments, allowing businesses to use augmented reality devices as a way to preview their products in the real world. Similarly, it can also be used to demo what products may look like in an environment for customers, as demonstrated by companies such as Mountain Equipment Co-op or Lowe's who use augmented reality to allow customers to preview what their products might look like at home through the use of 3D models. Augmented reality (AR) differs from virtual reality (VR) in the sense that in AR part of the surrounding environment is 'real' and AR is just adding layers of virtual objects to the real environment. On the other hand, in VR the surrounding environment is completely virtual and computer generated. A demonstration of how AR layers objects onto the real world can be seen with augmented reality games. WallaMe is an augmented reality game application that allows users to hide messages in real environments, utilizing geolocation technology in order to enable users to hide messages wherever they may wish in the world. Such applications have many uses in the world, including in activism and artistic expression. History 1901: L. Frank Baum, an author, first mentions the idea of an electronic display/spectacles that overlays data onto real life (in this case 'people'). It is named a 'character marker'. 1957–62: Morton Heilig, a cinematographer, creates and patents a simulator called Sensorama with visuals, sound, vibration, and smell. 1968: Ivan Sutherland creates the first head-mounted display that has graphics rendered by a computer. 1975: Myron Krueger creates Videoplace to allow users to interact with virtual objects. 1980: The research by Gavan Lintern of the University of Illinois is the first published work to show the value of a heads up display for teaching real-world flight skills. 1980: Steve Mann creates the first wearable computer, a computer vision system with text and graphical overlays on a photographically mediated scene. 1986: Within IBM, Ron Feigenblatt describes the most widely experienced form of AR today (viz. "magic window," e.g. smartphone-based Pokémon Go), use of a small, "smart" flat panel display positioned and oriented by hand. 1987: Douglas George and Robert Morris create a working prototype of an astronomical telescope-based "heads-up display" system (a precursor concept to augmented reality) which superimposed in the telescope eyepiece, over the actual sky images, multi-intensity star, and celestial body images, and other relevant information. 1990: The term augmented reality is attributed to Thomas P. Caudell, a former Boeing researcher. 1992: Louis Rosenberg developed one of the first functioning AR systems, called Virtual Fixtures, at the United States Air Force Research Laboratory—Armstrong, that demonstrated benefit to human perception. 1992: Steven Feiner, Blair MacIntyre and Doree Seligmann present an early paper on an AR system prototype, KARMA, at the Graphics Interface conference. 1993: The CMOS active-pixel sensor, a type of metal–oxide–semiconductor (MOS) image sensor, was developed at NASA's Jet Propulsion Laboratory. CMOS sensors are later widely used for optical tracking in AR technology. 1993: Mike Abernathy, et al., report the first use of augmented reality in identifying space debris using Rockwell WorldView by overlaying satellite geographic trajectories on live telescope video. 1993: A widely cited version of the paper above is published in Communications of the ACM – Special issue on computer augmented environments, edited by Pierre Wellner, Wendy Mackay, and Rich Gold. 1993: Loral WDL, with sponsorship from STRICOM, performed the first demonstration combining live AR-equipped vehicles and manned simulators. Unpublished paper, J. Barrilleaux, "Experiences and Observations in Applying Augmented Reality to Live Training", 1999. 1994: Julie Martin creates first 'Augmented Reality Theater production', Dancing in Cyberspace, funded by the Australia Council for the Arts, features dancers and acrobats manipulating body–sized virtual object in real time, projected into the same physical space and performance plane. The acrobats appeared immersed within the virtual object and environments. The installation used Silicon Graphics computers and Polhemus sensing system. 1996: General Electric develops system for projecting information from 3D CAD models onto real-world instances of those models. 1998: Spatial augmented reality introduced at University of North Carolina at Chapel Hill by Ramesh Raskar, Greg Welch, Henry Fuchs. 1999: Frank Delgado, Mike Abernathy et al. report successful flight test of LandForm software video map overlay from a helicopter at Army Yuma Proving Ground overlaying video with runways, taxiways, roads and road names. 1999: The US Naval Research Laboratory engages on a decade-long research program called the Battlefield Augmented Reality System (BARS) to prototype some of the early wearable systems for dismounted soldier operating in urban environment for situation awareness and training. 1999: NASA X-38 flown using LandForm software video map overlays at Dryden Flight Research Center. 2000: Rockwell International Science Center demonstrates tetherless wearable augmented reality systems receiving analog video and 3-D audio over radio-frequency wireless channels. The systems incorporate outdoor navigation capabilities, with digital horizon silhouettes from a terrain database overlain in real time on the live outdoor scene, allowing visualization of terrain made invisible by clouds and fog. 2004: An outdoor helmet-mounted AR system was demonstrated by Trimble Navigation and the Human Interface Technology Laboratory (HIT lab). 2006: Outland Research develops AR media player that overlays virtual content onto a users view of the real world synchronously with playing music, thereby providing an immersive AR entertainment experience. 2008: Wikitude AR Travel Guide launches on 20 Oct 2008 with the G1 Android phone. 2009: ARToolkit was ported to Adobe Flash (FLARToolkit) by Saqoosha, bringing augmented reality to the web browser. 2012: Launch of Lyteshot, an interactive AR gaming platform that utilizes smart glasses for game data 2015: Microsoft announced the HoloLens augmented reality headset, which uses various sensors and a processing unit to display virtual imagery over the real world. 2015: Snap, Inc. releases "Lenses", augmented reality filters in the Snapchat application. 2016: Niantic released Pokémon Go for iOS and Android in July 2016. The game quickly became one of the most popular smartphone applications and in turn spikes the popularity of augmented reality games. 2018: Magic Leap launched the Magic Leap One augmented reality headset. Leap Motion announced the Project North Star augmented reality headset, and later released it under an open source license. 2019: Microsoft announced HoloLens 2 with significant improvements in terms of field of view and ergonomics. 2022: Magic Leap launched the Magic Leap 2 headset. 2024: Meta Platforms revealed the Orion AR glasses prototype. Hardware Augmented reality requires hardware components including a processor, display, sensors, and input devices. Modern mobile computing devices like smartphones and tablet computers contain these elements, which often include a camera and microelectromechanical systems (MEMS) sensors such as an accelerometer, GPS, and solid state compass, making them suitable AR platforms. Displays Various technologies can be used to display augmented reality, including optical projection systems, monitors, and handheld devices. Two of the display technologies used in augmented reality are diffractive waveguides and reflective waveguides. A head-mounted display (HMD) is a display device worn on the forehead, such as a harness or helmet-mounted. HMDs place images of both the physical world and virtual objects over the user's field of view. Modern HMDs often employ sensors for six degrees of freedom monitoring that allow the system to align virtual information to the physical world and adjust accordingly with the user's head movements. When using AR technology, the HMDs only require relatively small displays. In this situation, liquid crystals on silicon (LCOS) and micro-OLED (organic light-emitting diodes) are commonly used. HMDs can provide VR users with mobile and collaborative experiences. Specific providers, such as uSens and Gestigon, include gesture controls for full virtual immersion. Vuzix is a company that has produced a number of head-worn optical see through displays marketed for augmented reality. Eyeglasses AR displays can be rendered on devices resembling eyeglasses. Versions include eyewear that employs cameras to intercept the real world view and re-display its augmented view through the eyepieces and devices in which the AR imagery is projected through or reflected off the surfaces of the eyewear lens pieces. The EyeTap (also known as Generation-2 Glass) captures rays of light that would otherwise pass through the center of the lens of the wearer's eye, and substitutes synthetic computer-controlled light for each ray of real light. The Generation-4 Glass (Laser EyeTap) is similar to the VRD (i.e. it uses a computer-controlled laser light source) except that it also has infinite depth of focus and causes the eye itself to, in effect, function as both a camera and a display by way of exact alignment with the eye and resynthesis (in laser light) of rays of light entering the eye. HUD A head-up display (HUD) is a transparent display that presents data without requiring users to look away from their usual viewpoints. A precursor technology to augmented reality, heads-up displays were first developed for pilots in the 1950s, projecting simple flight data into their line of sight, thereby enabling them to keep their "heads up" and not look down at the instruments. Near-eye augmented reality devices can be used as portable head-up displays as they can show data, information, and images while the user views the real world. Many definitions of augmented reality only define it as overlaying the information. This is basically what a head-up display does; however, practically speaking, augmented reality is expected to include registration and tracking between the superimposed perceptions, sensations, information, data, and images and some portion of the real world. Contact lenses Contact lenses that display AR imaging are in development. These bionic contact lenses might contain the elements for display embedded into the lens including integrated circuitry, LEDs and an antenna for wireless communication. The first contact lens display was patented in 1999 by Steve Mann and was intended to work in combination with AR spectacles, but the project was abandoned, then 11 years later in 2010–2011. Another version of contact lenses, in development for the U.S. military, is designed to function with AR spectacles, allowing soldiers to focus on close-to-the-eye AR images on the spectacles and distant real world objects at the same time. At CES 2013, a company called Innovega also unveiled similar contact lenses that required being combined with AR glasses to work. Many scientists have been working on contact lenses capable of different technological feats. A patent filed by Samsung describes an AR contact lens, that, when finished, will include a built-in camera on the lens itself. The design is intended to control its interface by blinking an eye. It is also intended to be linked with the user's smartphone to review footage, and control it separately. When successful, the lens would feature a camera, or sensor inside of it. It is said that it could be anything from a light sensor, to a temperature sensor. The first publicly unveiled working prototype of an AR contact lens not requiring the use of glasses in conjunction was developed by Mojo Vision and announced and shown off at CES 2020. Virtual retinal display A virtual retinal display (VRD) is a personal display device under development at the University of Washington's Human Interface Technology Laboratory under Dr. Thomas A. Furness III. With this technology, a display is scanned directly onto the retina of a viewer's eye. This results in bright images with high resolution and high contrast. The viewer sees what appears to be a conventional display floating in space. Several of tests were done to analyze the safety of the VRD. In one test, patients with partial loss of vision—having either macular degeneration (a disease that degenerates the retina) or keratoconus—were selected to view images using the technology. In the macular degeneration group, five out of eight subjects preferred the VRD images to the cathode-ray tube (CRT) or paper images and thought they were better and brighter and were able to see equal or better resolution levels. The Keratoconus patients could all resolve smaller lines in several line tests using the VRD as opposed to their own correction. They also found the VRD images to be easier to view and sharper. As a result of these several tests, virtual retinal display is considered safe technology. Virtual retinal display creates images that can be seen in ambient daylight and ambient room light. The VRD is considered a preferred candidate to use in a surgical display due to its combination of high resolution and high contrast and brightness. Additional tests show high potential for VRD to be used as a display technology for patients that have low vision. Handheld A Handheld display employs a small display that fits in a user's hand. All handheld AR solutions to date opt for video see-through. Initially handheld AR employed fiducial markers, and later GPS units and MEMS sensors such as digital compasses and six degrees of freedom accelerometer–gyroscope. Today simultaneous localization and mapping (SLAM) markerless trackers such as PTAM (parallel tracking and mapping) are starting to come into use. Handheld display AR promises to be the first commercial success for AR technologies. The two main advantages of handheld AR are the portable nature of handheld devices and the ubiquitous nature of camera phones. The disadvantages are the physical constraints of the user having to hold the handheld device out in front of them at all times, as well as the distorting effect of classically wide-angled mobile phone cameras when compared to the real world as viewed through the eye. Projection mapping Projection mapping augments real-world objects and scenes without the use of special displays such as monitors, head-mounted displays or hand-held devices. Projection mapping makes use of digital projectors to display graphical information onto physical objects. The key difference in projection mapping is that the display is separated from the users of the system. Since the displays are not associated with each user, projection mapping scales naturally up to groups of users, allowing for collocated collaboration between users. Examples include shader lamps, mobile projectors, virtual tables, and smart projectors. Shader lamps mimic and augment reality by projecting imagery onto neutral objects. This provides the opportunity to enhance the object's appearance with materials of a simple unit—a projector, camera, and sensor. Other applications include table and wall projections. Virtual showcases, which employ beam splitter mirrors together with multiple graphics displays, provide an interactive means of simultaneously engaging with the virtual and the real. A projection mapping system can display on any number of surfaces in an indoor setting at once. Projection mapping supports both a graphical visualization and passive haptic sensation for the end users. Users are able to touch physical objects in a process that provides passive haptic sensation. Tracking Modern mobile augmented-reality systems use one or more of the following motion tracking technologies: digital cameras and/or other optical sensors, accelerometers, GPS, gyroscopes, solid state compasses, radio-frequency identification (RFID). These technologies offer varying levels of accuracy and precision. These technologies are implemented in the ARKit API by Apple and ARCore API by Google to allow tracking for their respective mobile device platforms. Input devices Techniques include speech recognition systems that translate a user's spoken words into computer instructions, and gesture recognition systems that interpret a user's body movements by visual detection or from sensors embedded in a peripheral device such as a wand, stylus, pointer, glove or other body wear. Products which are trying to serve as a controller of AR headsets include Wave by Seebright Inc. and Nimble by Intugine Technologies. Computer Computers are responsible for graphics in augmented reality. For camera-based 3D tracking methods, a computer analyzes the sensed visual and other data to synthesize and position virtual objects. With the improvement of technology and computers, augmented reality is going to lead to a drastic change on ones perspective of the real world. Computers are improving at a very fast rate, leading to new ways to improve other technology. Computers are the core of augmented reality. The computer receives data from the sensors which determine the relative position of an objects' surface. This translates to an input to the computer which then outputs to the users by adding something that would otherwise not be there. The computer comprises memory and a processor. The computer takes the scanned environment then generates images or a video and puts it on the receiver for the observer to see. The fixed marks on an object's surface are stored in the memory of a computer. The computer also withdraws from its memory to present images realistically to the onlooker. Projector Projectors can also be used to display AR contents. The projector can throw a virtual object on a projection screen and the viewer can interact with this virtual object. Projection surfaces can be many objects such as walls or glass panes. Networking Mobile augmented reality applications are gaining popularity because of the wide adoption of mobile and especially wearable devices. However, they often rely on computationally intensive computer vision algorithms with extreme latency requirements. To compensate for the lack of computing power, offloading data processing to a distant machine is often desired. Computation offloading introduces new constraints in applications, especially in terms of latency and bandwidth. Although there are a plethora of real-time multimedia transport protocols, there is a need for support from network infrastructure as well. Software and algorithms A key measure of AR systems is how realistically they integrate virtual imagery with the real world. The software must derive real world coordinates, independent of camera, and camera images. That process is called image registration, and uses different methods of computer vision, mostly related to video tracking. Many computer vision methods of augmented reality are inherited from visual odometry. Usually those methods consist of two parts. The first stage is to detect interest points, fiducial markers or optical flow in the camera images. This step can use feature detection methods like corner detection, blob detection, edge detection or thresholding, and other image processing methods. The second stage restores a real world coordinate system from the data obtained in the first stage. Some methods assume objects with known geometry (or fiducial markers) are present in the scene. In some of those cases the scene 3D structure should be calculated beforehand. If part of the scene is unknown simultaneous localization and mapping (SLAM) can map relative positions. If no information about scene geometry is available, structure from motion methods like bundle adjustment are used. Mathematical methods used in the second stage include: projective (epipolar) geometry, geometric algebra, rotation representation with exponential map, kalman and particle filters, nonlinear optimization, robust statistics. In augmented reality, the distinction is made between two distinct modes of tracking, known as marker and markerless. Markers are visual cues which trigger the display of the virtual information. A piece of paper with some distinct geometries can be used. The camera recognizes the geometries by identifying specific points in the drawing. Markerless tracking, also called instant tracking, does not use markers. Instead, the user positions the object in the camera view preferably in a horizontal plane. It uses sensors in mobile devices to accurately detect the real-world environment, such as the locations of walls and points of intersection. Augmented Reality Markup Language (ARML) is a data standard developed within the Open Geospatial Consortium (OGC), which consists of Extensible Markup Language (XML) grammar to describe the location and appearance of virtual objects in the scene, as well as ECMAScript bindings to allow dynamic access to properties of virtual objects. To enable rapid development of augmented reality applications, software development applications have emerged, including Lens Studio from Snapchat and Spark AR from Facebook. Augmented reality Software Development Kits (SDKs) have been launched by Apple and Google. Development AR systems rely heavily on the immersion of the user. The following lists some considerations for designing augmented reality applications: Environmental/context design Context Design focuses on the end-user's physical surrounding, spatial space, and accessibility that may play a role when using the AR system. Designers should be aware of the possible physical scenarios the end-user may be in such as: Public, in which the users use their whole body to interact with the software Personal, in which the user uses a smartphone in a public space Intimate, in which the user is sitting with a desktop and is not really moving Private, in which the user has on a wearable. By evaluating each physical scenario, potential safety hazards can be avoided and changes can be made to greater improve the end-user's immersion. UX designers will have to define user journeys for the relevant physical scenarios and define how the interface reacts to each. Another aspect of context design involves the design of the system's functionality and its ability to accommodate user preferences. While accessibility tools are common in basic application design, some consideration should be made when designing time-limited prompts (to prevent unintentional operations), audio cues and overall engagement time. In some situations, the application's functionality may hinder the user's ability. For example, applications that is used for driving should reduce the amount of user interaction and use audio cues instead. Interaction design Interaction design in augmented reality technology centers on the user's engagement with the end product to improve the overall user experience and enjoyment. The purpose of interaction design is to avoid alienating or confusing the user by organizing the information presented. Since user interaction relies on the user's input, designers must make system controls easier to understand and accessible. A common technique to improve usability for augmented reality applications is by discovering the frequently accessed areas in the device's touch display and design the application to match those areas of control. It is also important to structure the user journey maps and the flow of information presented which reduce the system's overall cognitive load and greatly improves the learning curve of the application. In interaction design, it is important for developers to utilize augmented reality technology that complement the system's function or purpose. For instance, the utilization of exciting AR filters and the design of the unique sharing platform in Snapchat enables users to augment their in-app social interactions. In other applications that require users to understand the focus and intent, designers can employ a reticle or raycast from the device. Visual design To improve the graphic interface elements and user interaction, developers may use visual cues to inform the user what elements of UI are designed to interact with and how to interact with them. Visual cue design can make interactions seem more natural. In some augmented reality applications that use a 2D device as an interactive surface, the 2D control environment does not translate well in 3D space, which can make users hesitant to explore their surroundings. To solve this issue, designers should apply visual cues to assist and encourage users to explore their surroundings. It is important to note the two main objects in AR when developing VR applications: 3D volumetric objects that are manipulated and realistically interact with light and shadow; and animated media imagery such as images and videos which are mostly traditional 2D media rendered in a new context for augmented reality. When virtual objects are projected onto a real environment, it is challenging for augmented reality application designers to ensure a perfectly seamless integration relative to the real-world environment, especially with 2D objects. As such, designers can add weight to objects, use depths maps, and choose different material properties that highlight the object's presence in the real world. Another visual design that can be applied is using different lighting techniques or casting shadows to improve overall depth judgment. For instance, a common lighting technique is simply placing a light source overhead at the 12 o’clock position, to create shadows on virtual objects. Uses Augmented reality has been explored for many uses, including gaming, medicine, and entertainment. It has also been explored for education and business. Example application areas described below include archaeology, architecture, commerce and education. Some of the earliest cited examples include augmented reality used to support surgery by providing virtual overlays to guide medical practitioners, to AR content for astronomy and welding. Archaeology AR has been used to aid archaeological research. By augmenting archaeological features onto the modern landscape, AR allows archaeologists to formulate possible site configurations from extant structures. Computer generated models of ruins, buildings, landscapes or even ancient people have been recycled into early archaeological AR applications. For example, implementing a system like VITA (Visual Interaction Tool for Archaeology) will allow users to imagine and investigate instant excavation results without leaving their home. Each user can collaborate by mutually "navigating, searching, and viewing data". Hrvoje Benko, a researcher in the computer science department at Columbia University, points out that these particular systems and others like them can provide "3D panoramic images and 3D models of the site itself at different excavation stages" all the while organizing much of the data in a collaborative way that is easy to use. Collaborative AR systems supply multimodal interactions that combine the real world with virtual images of both environments. Architecture AR can aid in visualizing building projects. Computer-generated images of a structure can be superimposed onto a real-life local view of a property before the physical building is constructed there; this was demonstrated publicly by Trimble Navigation in 2004. AR can also be employed within an architect's workspace, rendering animated 3D visualizations of their 2D drawings. Architecture sight-seeing can be enhanced with AR applications, allowing users viewing a building's exterior to virtually see through its walls, viewing its interior objects and layout. With continual improvements to GPS accuracy, businesses are able to use augmented reality to visualize georeferenced models of construction sites, underground structures, cables and pipes using mobile devices. Augmented reality is applied to present new projects, to solve on-site construction challenges, and to enhance promotional materials. Examples include the Daqri Smart Helmet, an Android-powered hard hat used to create augmented reality for the industrial worker, including visual instructions, real-time alerts, and 3D mapping. Following the Christchurch earthquake, the University of Canterbury released CityViewAR, which enabled city planners and engineers to visualize buildings that had been destroyed. This not only provided planners with tools to reference the previous cityscape, but it also served as a reminder of the magnitude of the resulting devastation, as entire buildings had been demolished. Education and Training In educational settings, AR has been used to complement a standard curriculum. Text, graphics, video, and audio may be superimposed into a student's real-time environment. Textbooks, flashcards and other educational reading material may contain embedded "markers" or triggers that, when scanned by an AR device, produced supplementary information to the student rendered in a multimedia format. The 2015 Virtual, Augmented and Mixed Reality: 7th International Conference mentioned Google Glass as an example of augmented reality that can replace the physical classroom. First, AR technologies help learners engage in authentic exploration in the real world, and virtual objects such as texts, videos, and pictures are supplementary elements for learners to conduct investigations of the real-world surroundings. As AR evolves, students can participate interactively and interact with knowledge more authentically. Instead of remaining passive recipients, students can become active learners, able to interact with their learning environment. Computer-generated simulations of historical events allow students to explore and learning details of each significant area of the event site. In higher education, Construct3D, a Studierstube system, allows students to learn mechanical engineering concepts, math or geometry. Chemistry AR apps allow students to visualize and interact with the spatial structure of a molecule using a marker object held in the hand. Others have used HP Reveal, a free app, to create AR notecards for studying organic chemistry mechanisms or to create virtual demonstrations of how to use laboratory instrumentation. Anatomy students can visualize different systems of the human body in three dimensions. Using AR as a tool to learn anatomical structures has been shown to increase the learner knowledge and provide intrinsic benefits, such as increased engagement and learner immersion. AR has been used to develop different safety training applications for several types of disasters, such as, earthquakes and building fire, and health and safety tasks. Further, several AR solutions have been proposed and tested to navigate building evacuees towards safe places in both large scale and small scale disasters. AR applications can have several overlapping with many other digital technologies, such as BIM, internet of things and artificial intelligence, to generate smarter safety training and navigation solutions. Industrial manufacturing AR is used to substitute paper manuals with digital instructions which are overlaid on the manufacturing operator's field of view, reducing mental effort required to operate. AR makes machine maintenance efficient because it gives operators direct access to a machine's maintenance history. Virtual manuals help manufacturers adapt to rapidly-changing product designs, as digital instructions are more easily edited and distributed compared to physical manuals. Digital instructions increase operator safety by removing the need for operators to look at a screen or manual away from the working area, which can be hazardous. Instead, the instructions are overlaid on the working area. The use of AR can increase operators' feeling of safety when working near high-load industrial machinery by giving operators additional information on a machine's status and safety functions, as well as hazardous areas of the workspace. Commerce AR is used to integrate print and video marketing. Printed marketing material can be designed with certain "trigger" images that, when scanned by an AR-enabled device using image recognition, activate a video version of the promotional material. A major difference between augmented reality and straightforward image recognition is that one can overlay multiple media at the same time in the view screen, such as social media share buttons, the in-page video even audio and 3D objects. Traditional print-only publications are using augmented reality to connect different types of media. AR can enhance product previews such as allowing a customer to view what's inside a product's packaging without opening it. AR can also be used as an aid in selecting products from a catalog or through a kiosk. Scanned images of products can activate views of additional content such as customization options and additional images of the product in its use. By 2010, virtual dressing rooms had been developed for e-commerce. In 2012, a mint used AR techniques to market a commemorative coin for Aruba. The coin itself was used as an AR trigger, and when held in front of an AR-enabled device it revealed additional objects and layers of information that were not visible without the device. In 2018, Apple announced Universal Scene Description (USDZ) AR file support for iPhones and iPads with iOS 12. Apple has created an AR QuickLook Gallery that allows masses to experience augmented reality on their own Apple device. In 2018, Shopify, the Canadian e-commerce company, announced AR Quick Look integration. Their merchants will be able to upload 3D models of their products and their users will be able to tap on the models inside the Safari browser on their iOS devices to view them in their real-world environments. In 2018, Twinkl released a free AR classroom application. Pupils can see how York looked over 1,900 years ago. Twinkl launched the first ever multi-player AR game, Little Red and has over 100 free AR educational models. Augmented reality is becoming more frequently used for online advertising. Retailers offer the ability to upload a picture on their website and "try on" various clothes which are overlaid on the picture. Even further, companies such as Bodymetrics install dressing booths in department stores that offer full-body scanning. These booths render a 3-D model of the user, allowing the consumers to view different outfits on themselves without the need of physically changing clothes. For example, JC Penney and Bloomingdale's use "virtual dressing rooms" that allow customers to see themselves in clothes without trying them on. Another store that uses AR to market clothing to its customers is Neiman Marcus. Neiman Marcus offers consumers the ability to see their outfits in a 360-degree view with their "memory mirror". Makeup stores like L'Oreal, Sephora, Charlotte Tilbury, and Rimmel also have apps that utilize AR. These apps allow consumers to see how the makeup will look on them. According to Greg Jones, director of AR and VR at Google, augmented reality is going to "reconnect physical and digital retail". AR technology is also used by furniture retailers such as IKEA, Houzz, and Wayfair. These retailers offer apps that allow consumers to view their products in their home prior to purchasing anything. In 2017, Ikea announced the Ikea Place app. It contains a catalogue of over 2,000 products—nearly the company's full collection of sofas, armchairs, coffee tables, and storage units which one can place anywhere in a room with their phone. The app made it possible to have 3D and true-to-scale models of furniture in the customer's living space. IKEA realized that their customers are not shopping in stores as often or making direct purchases anymore. Shopify's acquisition of Primer, an AR app aims to push small and medium-sized sellers towards interactive AR shopping with easy to use AR integration and user experience for both merchants and consumers. AR helps the retail industry reduce operating costs. Merchants upload product information to the AR system, and consumers can use mobile terminals to search and generate 3D maps. Literature The first description of AR as it is known today was in Virtual Light, the 1994 novel by William Gibson. In 2011, AR was blended with poetry by ni ka from Sekai Camera in Tokyo, Japan. The prose of these AR poems come from Paul Celan, Die Niemandsrose, expressing the aftermath of the 2011 Tōhoku earthquake and tsunami. Visual art AR applied in the visual arts allows objects or places to trigger artistic multidimensional experiences and interpretations of reality. The Australian new media artist Jeffrey Shaw pioneered Augmented Reality in three artworks: Viewpoint in 1975, Virtual Sculptures in 1987 and The Golden Calf in 1993. He continues to explore new permutations of AR in numerous recent works. Manifest.AR was an international artists' collective founded in 2010 that specialized in augmented reality (AR) art and interventions. The collective typically created site-specific AR installations that could be viewed through mobile devices using custom-developed applications. Their work often challenged traditional notions of art exhibition and ownership by placing virtual artworks in spaces without institutional permission. The collective gained prominence in 2010 when they staged an unauthorized virtual exhibition at the Museum of Modern Art (MoMA) in New York City, overlaying their digital artworks throughout the museum's spaces using AR technology. The collective's unauthorized AR intervention at MoMA involved placing virtual artworks throughout the museum's spaces, viewable through mobile devices. In 2011, members of Manifest.AR created AR artworks that were virtually placed throughout the Venice Biennial, creating an unofficial parallel exhibition accessible through mobile devices. During the Occupy Wall Street movement in 2011, the collective created AR installations in and around Zuccotti Park, adding a digital dimension to the physical protests. Key members of the collective have included: Mark Skwarek; John Craig Freeman; Will Pappenheimer; Tamiko Thiel; and Sander Veenhof. The group published their "AR Art Manifesto" in 2011, which outlined their artistic philosophy and approach to augmented reality as a medium. The manifesto emphasized the democratic potential of AR technology and its ability to challenge traditional institutional control over public space and art display. Manifest.AR has been influential in: Pioneering artistic applications of AR technology; Developing new forms of institutional critique; Expanding concepts of public art and digital space; and Influencing subsequent generations of new media artists. Their work has been documented and discussed in various publications about digital art and new media, and has influenced contemporary discussions about virtual and augmented reality in artistic practice. Augmented reality can aid in the progression of visual art in museums by allowing museum visitors to view artwork in galleries in a multidimensional way through their phone screens. The Museum of Modern Art in New York has created an exhibit in their art museum showcasing AR features that viewers can see using an app on their smartphone. The museum has developed their personal app, called MoMAR Gallery, that museum guests can download and use in the augmented reality specialized gallery in order to view the museum's paintings in a different way. This allows individuals to see hidden aspects and information about the paintings, and to be able to have an interactive technological experience with artwork as well. AR technology was used in Nancy Baker Cahill's "Margin of Error" and "Revolutions," the two public art pieces she created for the 2019 Desert X exhibition. AR technology aided the development of eye tracking technology to translate a disabled person's eye movements into drawings on a screen. A Danish artist, Olafur Eliasson, has placed objects like burning suns, extraterrestrial rocks, and rare animals, into the user's environment. Martin & Muñoz started using Augmented Reality (AR) technology in 2020 to create and place virtual works, based on their snow globes, in their exhibitions and in user's environments. Their first AR work was presented at the Cervantes Institute in New York in early 2022. Fitness AR hardware and software for use in fitness includes smart glasses made for biking and running, with performance analytics and map navigation projected onto the user's field of vision, and boxing, martial arts, and tennis, where users remain aware of their physical environment for safety. Fitness-related games and software include Pokémon Go and Jurassic World Alive. Human–computer interaction Human–computer interaction (HCI) is an interdisciplinary area of computing that deals with design and implementation of systems that interact with people. Researchers in HCI come from a number of disciplines, including computer science, engineering, design, human factor, and social science, with a shared goal to solve problems in the design and the use of technology so that it can be used more easily, effectively, efficiently, safely, and with satisfaction. According to a 2017 Time article, in about 15 to 20 years it is predicted that augmented reality and virtual reality are going to become the primary use for computer interactions. Remote collaboration Primary school children learn easily from interactive experiences. As an example, astronomical constellations and the movements of objects in the solar system were oriented in 3D and overlaid in the direction the device was held, and expanded with supplemental video information. Paper-based science book illustrations could seem to come alive as video without requiring the child to navigate to web-based materials. In 2013, a project was launched on Kickstarter to teach about electronics with an educational toy that allowed children to scan their circuit with an iPad and see the electric current flowing around. While some educational apps were available for AR by 2016, it was not broadly used. Apps that leverage augmented reality to aid learning included SkyView for studying astronomy, AR Circuits for building simple electric circuits, and SketchAR for drawing. AR would also be a way for parents and teachers to achieve their goals for modern education, which might include providing more individualized and flexible learning, making closer connections between what is taught at school and the real world, and helping students to become more engaged in their own learning. Emergency management/search and rescue Augmented reality systems are used in public safety situations, from super storms to suspects at large. As early as 2009, two articles from Emergency Management discussed AR technology for emergency management. The first was "Augmented Reality—Emerging Technology for Emergency Management", by Gerald Baron. According to Adam Crow,: "Technologies like augmented reality (ex: Google Glass) and the growing expectation of the public will continue to force professional emergency managers to radically shift when, where, and how technology is deployed before, during, and after disasters." Another early example was a search aircraft looking for a lost hiker in rugged mountain terrain. Augmented reality systems provided aerial camera operators with a geographic awareness of forest road names and locations blended with the camera video. The camera operator was better able to search for the hiker knowing the geographic context of the camera image. Once located, the operator could more efficiently direct rescuers to the hiker's location because the geographic position and reference landmarks were clearly labeled. Social interaction AR can be used to facilitate social interaction. An augmented reality social network framework called Talk2Me enables people to disseminate information and view others' advertised information in an augmented reality way. The timely and dynamic information sharing and viewing functionalities of Talk2Me help initiate conversations and make friends for users with people in physical proximity. However, use of an AR headset can inhibit the quality of an interaction between two people if one isn't wearing one if the headset becomes a distraction. Augmented reality also gives users the ability to practice different forms of social interactions with other people in a safe, risk-free environment. Hannes Kauffman, Associate Professor for virtual reality at TU Vienna, says: "In collaborative augmented reality multiple users may access a shared space populated by virtual objects, while remaining grounded in the real world. This technique is particularly powerful for educational purposes when users are collocated and can use natural means of communication (speech, gestures, etc.), but can also be mixed successfully with immersive VR or remote collaboration." Hannes cites education as a potential use of this technology. Video games The gaming industry embraced AR technology. A number of games were developed for prepared indoor environments, such as AR air hockey, Titans of Space, collaborative combat against virtual enemies, and AR-enhanced pool table games. In 2010, Ogmento became the first AR gaming startup to receive VC Funding. The company went on to produce early location-based AR games for titles like Paranormal Activity: Sanctuary, NBA: King of the Court, and Halo: King of the Hill. The companies computer vision technology was eventually repackaged and sold to Apple, became a major contribution to ARKit. Augmented reality allows video game players to experience digital game play in a real-world environment. Niantic released the augmented reality mobile game Pokémon Go. Disney has partnered with Lenovo to create the augmented reality game Star Wars: Jedi Challenges that works with a Lenovo Mirage AR headset, a tracking sensor and a Lightsaber controller, scheduled to launch in December 2017. Industrial design AR allows industrial designers to experience a product's design and operation before completion. Volkswagen has used AR for comparing calculated and actual crash test imagery. AR has been used to visualize and modify car body structure and engine layout. It has also been used to compare digital mock-ups with physical mock-ups to find discrepancies between them. Healthcare planning, practice and education One of the first applications of augmented reality was in healthcare, particularly to support the planning, practice, and training of surgical procedures. As far back as 1992, enhancing human performance during surgery was a formally stated objective when building the first augmented reality systems at U.S. Air Force laboratories. Since 2005, a device called a near-infrared vein finder that films subcutaneous veins, processes and projects the image of the veins onto the skin has been used to locate veins. AR provides surgeons with patient monitoring data in the style of a fighter pilot's heads-up display, and allows patient imaging records, including functional videos, to be accessed and overlaid. Examples include a virtual X-ray view based on prior tomography or on real-time images from ultrasound and confocal microscopy probes, visualizing the position of a tumor in the video of an endoscope, or radiation exposure risks from X-ray imaging devices. AR can enhance viewing a fetus inside a mother's womb. Siemens, Karl Storz and IRCAD have developed a system for laparoscopic liver surgery that uses AR to view sub-surface tumors and vessels. AR has been used for cockroach phobia treatment and to reduce the fear of spiders. Patients wearing augmented reality glasses can be reminded to take medications. Augmented reality can be very helpful in the medical field. It could be used to provide crucial information to a doctor or surgeon without having them take their eyes off the patient. On 30 April 2015 Microsoft announced the Microsoft HoloLens, their first attempt at augmented reality. The HoloLens has advanced through the years and is capable of projecting holograms for near infrared fluorescence based image guided surgery. As augmented reality advances, it finds increasing applications in healthcare. Augmented reality and similar computer based-utilities are being used to train medical professionals. In healthcare, AR can be used to provide guidance during diagnostic and therapeutic interventions e.g. during surgery. Magee et al., for instance, describe the use of augmented reality for medical training in simulating ultrasound-guided needle placement. Similarly, Javaid, Mohd, Haleem, and Abid found that virtual reality provided medical students' brains with an experience that simulates motion and the surgery experience. A very recent study by Akçayır, Akçayır, Pektaş, and Ocak (2016) revealed that AR technology both improves university students' laboratory skills and helps them to build positive attitudes relating to physics laboratory work. Recently, augmented reality began seeing adoption in neurosurgery, a field that requires heavy amounts of imaging before procedures. Spatial immersion and interaction Augmented reality applications, running on handheld devices utilized as virtual reality headsets, can also digitize human presence in space and provide a computer generated model of them, in a virtual space where they can interact and perform various actions. Such capabilities are demonstrated by Project Anywhere, developed by a postgraduate student at ETH Zurich, which was dubbed as an "out-of-body experience". Flight training Building on decades of perceptual-motor research in experimental psychology, researchers at the Aviation Research Laboratory of the University of Illinois at Urbana–Champaign used augmented reality in the form of a flight path in the sky to teach flight students how to land an airplane using a flight simulator. An adaptive augmented schedule in which students were shown the augmentation only when they departed from the flight path proved to be a more effective training intervention than a constant schedule. Flight students taught to land in the simulator with the adaptive augmentation learned to land a light aircraft more quickly than students with the same amount of landing training in the simulator but with constant augmentation or without any augmentation. Military An interesting early application of AR occurred when Rockwell International created video map overlays of satellite and orbital debris tracks to aid in space observations at Air Force Maui Optical System. In their 1993 paper "Debris Correlation Using the Rockwell WorldView System" the authors describe the use of map overlays applied to video from space surveillance telescopes. The map overlays indicated the trajectories of various objects in geographic coordinates. This allowed telescope operators to identify satellites, and also to identify and catalog potentially dangerous space debris. Starting in 2003 the US Army integrated the SmartCam3D augmented reality system into the Shadow Unmanned Aerial System to aid sensor operators using telescopic cameras to locate people or points of interest. The system combined fixed geographic information including street names, points of interest, airports, and railroads with live video from the camera system. The system offered a "picture in picture" mode that allows it to show a synthetic view of the area surrounding the camera's field of view. This helps solve a problem in which the field of view is so narrow that it excludes important context, as if "looking through a soda straw". The system displays real-time friend/foe/neutral location markers blended with live video, providing the operator with improved situational awareness. Researchers at USAF Research Lab (Calhoun, Draper et al.) found an approximately two-fold increase in the speed at which UAV sensor operators found points of interest using this technology. This ability to maintain geographic awareness quantitatively enhances mission efficiency. The system is in use on the US Army RQ-7 Shadow and the MQ-1C Gray Eagle Unmanned Aerial Systems. In combat, AR can serve as a networked communication system that renders useful battlefield data onto a soldier's goggles in real time. From the soldier's viewpoint, people and various objects can be marked with special indicators to warn of potential dangers. Virtual maps and 360° view camera imaging can also be rendered to aid a soldier's navigation and battlefield perspective, and this can be transmitted to military leaders at a remote command center. The combination of 360° view cameras visualization and AR can be used on board combat vehicles and tanks as circular review system. AR can be an effective tool for virtually mapping out the 3D topologies of munition storages in the terrain, with the choice of the munitions combination in stacks and distances between them with a visualization of risk areas. The scope of AR applications also includes visualization of data from embedded munitions monitoring sensors. Navigation The NASA X-38 was flown using a hybrid synthetic vision system that overlaid map data on video to provide enhanced navigation for the spacecraft during flight tests from 1998 to 2002. It used the LandForm software which was useful for times of limited visibility, including an instance when the video camera window frosted over leaving astronauts to rely on the map overlays. The LandForm software was also test flown at the Army Yuma Proving Ground in 1999. In the photo at right one can see the map markers indicating runways, air traffic control tower, taxiways, and hangars overlaid on the video. AR can augment the effectiveness of navigation devices. Information can be displayed on an automobile's windshield indicating destination directions and meter, weather, terrain, road conditions and traffic information as well as alerts to potential hazards in their path. Since 2012, a Swiss-based company WayRay has been developing holographic AR navigation systems that use holographic optical elements for projecting all route-related information including directions, important notifications, and points of interest right into the drivers' line of sight and far ahead of the vehicle. Aboard maritime vessels, AR can allow bridge watch-standers to continuously monitor important information such as a ship's heading and speed while moving throughout the bridge or performing other tasks. Workplace Augmented reality may have a positive impact on work collaboration as people may be inclined to interact more actively with their learning environment. It may also encourage tacit knowledge renewal which makes firms more competitive. AR was used to facilitate collaboration among distributed team members via conferences with local and virtual participants. AR tasks included brainstorming and discussion meetings utilizing common visualization via touch screen tables, interactive digital whiteboards, shared design spaces and distributed control rooms. In industrial environments, augmented reality is proving to have a substantial impact with more and more use cases emerging across all aspect of the product lifecycle, starting from product design and new product introduction (NPI) to manufacturing to service and maintenance, to material handling and distribution. For example, labels were displayed on parts of a system to clarify operating instructions for a mechanic performing maintenance on a system. Assembly lines benefited from the usage of AR. In addition to Boeing, BMW and Volkswagen were known for incorporating this technology into assembly lines for monitoring process improvements. Big machines are difficult to maintain because of their multiple layers or structures. AR permits people to look through the machine as if with an x-ray, pointing them to the problem right away. As AR technology has evolved and second and third generation AR devices come to market, the impact of AR in enterprise continues to flourish. In the Harvard Business Review, Magid Abraham and Marco Annunziata discuss how AR devices are now being used to "boost workers' productivity on an array of tasks the first time they're used, even without prior training". They contend that "these technologies increase productivity by making workers more skilled and efficient, and thus have the potential to yield both more economic growth and better jobs". Broadcast and live events Weather visualizations were the first application of augmented reality in television. It has now become common in weather casting to display full motion video of images captured in real-time from multiple cameras and other imaging devices. Coupled with 3D graphics symbols and mapped to a common virtual geospatial model, these animated visualizations constitute the first true application of AR to TV. AR has become common in sports telecasting. Sports and entertainment venues are provided with see-through and overlay augmentation through tracked camera feeds for enhanced viewing by the audience. Examples include the yellow "first down" line seen in television broadcasts of American football games showing the line the offensive team must cross to receive a first down. AR is also used in association with football and other sporting events to show commercial advertisements overlaid onto the view of the playing area. Sections of rugby fields and cricket pitches also display sponsored images. Swimming telecasts often add a line across the lanes to indicate the position of the current record holder as a race proceeds to allow viewers to compare the current race to the best performance. Other examples include hockey puck tracking and annotations of racing car performance and snooker ball trajectories. AR has been used to enhance concert and theater performances. For example, artists allow listeners to augment their listening experience by adding their performance to that of other bands/groups of users. Tourism and sightseeing Travelers may use AR to access real-time informational displays regarding a location, its features, and comments or content provided by previous visitors. Advanced AR applications include simulations of historical events, places, and objects rendered into the landscape. AR applications linked to geographic locations present location information by audio, announcing features of interest at a particular site as they become visible to the user. Translation AR systems such as Word Lens can interpret the foreign text on signs and menus and, in a user's augmented view, re-display the text in the user's language. Spoken words of a foreign language can be translated and displayed in a user's view as printed subtitles. Music It has been suggested that augmented reality may be used in new methods of music production, mixing, control and visualization. In a proof-of-concept project Ian Sterling, an interaction design student at California College of the Arts, and software engineer Swaroop Pal demonstrated a HoloLens app whose primary purpose is to provide a 3D spatial UI for cross-platform devices—the Android Music Player app and Arduino-controlled Fan and Light—and also allow interaction using gaze and gesture control. Research by members of the CRIStAL at the University of Lille makes use of augmented reality to enrich musical performance. The ControllAR project allows musicians to augment their MIDI control surfaces with the remixed graphical user interfaces of music software. The Rouages project proposes to augment digital musical instruments to reveal their mechanisms to the audience and thus improve the perceived liveness. Reflets is a novel augmented reality display dedicated to musical performances where the audience acts as a 3D display by revealing virtual content on stage, which can also be used for 3D musical interaction and collaboration. Snapchat Snapchat users have access to augmented reality in the app through use of camera filters. In September 2017, Snapchat updated its app to include a camera filter that allowed users to render an animated, cartoon version of themselves called "Bitmoji". These animated avatars would be projected in the real world through the camera, and can be photographed or video recorded. In the same month, Snapchat also announced a new feature called "Sky Filters" that will be available on its app. This new feature makes use of augmented reality to alter the look of a picture taken of the sky, much like how users can apply the app's filters to other pictures. Users can choose from sky filters such as starry night, stormy clouds, beautiful sunsets, and rainbow. Concerns Reality modifications In a paper titled "Death by Pokémon GO", researchers at Purdue University's Krannert School of Management claim the game caused "a disproportionate increase in vehicular crashes and associated vehicular damage, personal injuries, and fatalities in the vicinity of locations, called PokéStops, where users can play the game while driving." Using data from one municipality, the paper extrapolates what that might mean nationwide and concluded "the increase in crashes attributable to the introduction of Pokémon GO is 145,632 with an associated increase in the number of injuries of 29,370 and an associated increase in the number of fatalities of 256 over the period of 6 July 2016, through 30 November 2016." The authors extrapolated the cost of those crashes and fatalities at between $2bn and $7.3 billion for the same period. Furthermore, more than one in three surveyed advanced Internet users would like to edit out disturbing elements around them, such as garbage or graffiti. They would like to even modify their surroundings by erasing street signs, billboard ads, and uninteresting shopping windows. So it seems that AR is as much a threat to companies as it is an opportunity. Although, this could be a nightmare to numerous brands that do not manage to capture consumer imaginations it also creates the risk that the wearers of augmented reality glasses may become unaware of surrounding dangers. Consumers want to use augmented reality glasses to change their surroundings into something that reflects their own personal opinions. Around two in five want to change the way their surroundings look and even how people appear to them. Next, to the possible privacy issues that are described below, overload and over-reliance issues are the biggest danger of AR. For the development of new AR-related products, this implies that the user-interface should follow certain guidelines as not to overload the user with information while also preventing the user from over-relying on the AR system such that important cues from the environment are missed. This is called the virtually-augmented key. Once the key is ignored, people might not desire the real world anymore. Privacy concerns The concept of modern augmented reality depends on the ability of the device to record and analyze the environment in real time. Because of this, there are potential legal concerns over privacy. While the First Amendment to the United States Constitution allows for such recording in the name of public interest, the constant recording of an AR device makes it difficult to do so without also recording outside of the public domain. Legal complications would be found in areas where a right to a certain amount of privacy is expected or where copyrighted media are displayed. In terms of individual privacy, there exists the ease of access to information that one should not readily possess about a given person. This is accomplished through facial recognition technology. Assuming that AR automatically passes information about persons that the user sees, there could be anything seen from social media, criminal record, and marital status. The Code of Ethics on Human Augmentation, which was originally introduced by Steve Mann in 2004 and further refined with Ray Kurzweil and Marvin Minsky in 2013, was ultimately ratified at the virtual reality Toronto conference on 25 June 2017. Property law The interaction of location-bound augmented reality with property law is largely undefined. Several models have been analysed for how this interaction may be resolved in a common law context: an extension of real property rights to also cover augmentations on or near the property with a strong notion of trespassing, forbidding augmentations unless allowed by the owner; an 'open range' system, where augmentations are allowed unless forbidden by the owner; and a 'freedom to roam' system, where real property owners have no control over non-disruptive augmentations. One issue experienced during the Pokémon Go craze was the game's players disturbing owners of private property while visiting nearby location-bound augmentations, which may have been on the properties or the properties may have been en route. The terms of service of Pokémon Go explicitly disclaim responsibility for players' actions, which may limit (but may not totally extinguish) the liability of its producer, Niantic, in the event of a player trespassing while playing the game: by Niantic's argument, the player is the one committing the trespass, while Niantic has merely engaged in permissible free speech. A theory advanced in lawsuits brought against Niantic is that their placement of game elements in places that will lead to trespass or an exceptionally large flux of visitors can constitute nuisance, despite each individual trespass or visit only being tenuously caused by Niantic. Another claim raised against Niantic is that the placement of profitable game elements on land without permission of the land's owners is unjust enrichment. More hypothetically, a property may be augmented with advertising or disagreeable content against its owner's wishes. Under American law, these situations are unlikely to be seen as a violation of real property rights by courts without an expansion of those rights to include augmented reality (similarly to how English common law came to recognise air rights). An article in the Michigan Telecommunications and Technology Law Review argues that there are three bases for this extension, starting with various understanding of property. The personality theory of property, outlined by Margaret Radin, is claimed to support extending property rights due to the intimate connection between personhood and ownership of property; however, her viewpoint is not universally shared by legal theorists. Under the utilitarian theory of property, the benefits from avoiding the harms to real property owners caused by augmentations and the tragedy of the commons, and the reduction in transaction costs by making discovery of ownership easy, were assessed as justifying recognising real property rights as covering location-bound augmentations, though there does remain the possibility of a tragedy of the anticommons from having to negotiate with property owners slowing innovation. Finally, following the 'property as the law of things' identification as supported by Thomas Merrill and Henry E Smith, location-based augmentation is naturally identified as a 'thing', and, while the non-rivalrous and ephemeral nature of digital objects presents difficulties to the excludeability prong of the definition, the article argues that this is not insurmountable. Some attempts at legislative regulation have been made in the United States. Milwaukee County, Wisconsin attempted to regulate augmented reality games played in its parks, requiring prior issuance of a permit, but this was criticised on free speech grounds by a federal judge; and Illinois considered mandating a notice and take down procedure for location-bound augmentations. An article for the Iowa Law Review observed that dealing with many local permitting processes would be arduous for a large-scale service, and, while the proposed Illinois mechanism could be made workable, it was reactive and required property owners to potentially continually deal with new augmented reality services; instead, a national-level geofencing registry, analogous to a do-not-call list, was proposed as the most desirable form of regulation to efficiently balance the interests of both providers of augmented reality services and real property owners. An article in the Vanderbilt Journal of Entertainment and Technology Law, however, analyses a monolithic do-not-locate registry as an insufficiently flexible tool, either permitting unwanted augmentations or foreclosing useful applications of augmented reality. Instead, it argues that an 'open range' model, where augmentations are permitted by default but property owners may restrict them on a case-by-case basis (and with noncompliance treated as a form of trespass), will produce the socially-best outcome. Notable researchers Ronald Azuma is a scientist and author of works on AR. Jeri Ellsworth headed a research effort for Valve on augmented reality (AR), later taking that research to her own start-up CastAR. The company, founded in 2013, eventually shuttered. Later, she created another start-up based on the same technology called Tilt Five; another AR start-up formed by her with the purpose of creating a device for digital board games. Steve Mann formulated an earlier concept of mediated reality in the 1970s and 1980s, using cameras, processors, and display systems to modify visual reality to help people see better (dynamic range management), building computerized welding helmets, as well as "augmediated reality" vision systems for use in everyday life. He is also an adviser to Meta. Dieter Schmalstieg and Daniel Wagner developed a marker tracking systems for mobile phones and PDAs in 2009. Ivan Sutherland invented the first VR head-mounted display at Harvard University. In media The futuristic short film Sight features contact lens-like augmented reality devices.	Technology	Computer science	null
13997033	https://en.wikipedia.org/wiki/Accretionary%20wedge	Accretionary wedge	An accretionary wedge or accretionary prism forms from sediments accreted onto the non-subducting tectonic plate at a convergent plate boundary. Most of the material in the accretionary wedge consists of marine sediments scraped off from the downgoing slab of oceanic crust, but in some cases the wedge includes the erosional products of volcanic island arcs formed on the overriding plate. An accretionary complex is a current (in modern use) or former accretionary wedge. Accretionary complexes are typically made up of a mix of turbidites of terrestrial material, basalts from the ocean floor, and pelagic and hemipelagic sediments. For example, most of the geological basement of Japan is made up of accretionary complexes. Materials within an accretionary wedge Accretionary wedges and accreted terranes are not equivalent to tectonic plates, but rather are associated with tectonic plates and accrete as a result of tectonic collision. Materials incorporated in accretionary wedges include: Ocean-floor basalts – typically seamounts scraped off the subducting plate Pelagic sediments – typically immediately overlying oceanic crust of the subducting plate Trench sediments – typically turbidites that may be derived from: Oceanic, volcanic island arc Continental volcanic arc and cordilleran orogen Adjacent continental masses located along strike (such as Barbados). Material transported into the trench by gravity sliding and debris flow from the forearc ridge (olistostrome) Piggy-back basins, which are small basins located in surface depression on the accretionary prism. Material exposed in the forearc ridge may include fragments of oceanic crust or high pressure metamorphic rocks thrust from deeper in the subduction zone. Elevated regions within the ocean basins such as linear island chains, ocean ridges, and small crustal fragments (such as Madagascar or Japan), known as terranes, are transported toward the subduction zone and accreted to the continental margin. Since the Late Devonian and Early Carboniferous periods, some 360 million years ago, subduction beneath the western margin of North America has resulted in several collisions with terranes, each producing a mountain-building event. The piecemeal addition of these accreted terranes has added an average of in width along the western margin of the North American continent. Geometry The topographic expression of the accretionary wedge forms a lip, which may dam basins of accumulated materials that, otherwise, would be transported into the trench from the overriding plate. Accretionary wedges are the home of mélange, intensely deformed packages of rocks that lack coherent internal layering and coherent internal order. The internal structure of an accretionary wedge is similar to that found in a thin-skinned foreland thrust belt. A series of thrusts verging towards the trench are formed with the youngest most outboard structures progressively uplifting the older more inboard thrusts. The shape of the wedge is determined by how readily the wedge will fail along its basal decollement and in its interior; this is highly sensitive to pore fluid pressure. This failure will result in a mature wedge that has an equilibrium triangular cross-sectional shape of a critical taper. Once the wedge reaches a critical taper, it will maintain that geometry and grow only into a larger similar triangle. Significance The small sections of oceanic crust that are thrust over the overriding plate are said to be obducted. Where this occurs, rare slices of ocean crust, known as ophiolites, are preserved on land. They provide a valuable natural laboratory for studying the composition and character of the oceanic crust and the mechanisms of their emplacement and preservation on land. A classic example is the Coast Range ophiolite of California, which is one of the most extensive ophiolite terranes in North America. This oceanic crust likely formed during the middle Jurassic Period, roughly 170 million years ago, in an extensional regime within either a back-arc or a forearc basin. It was later accreted to the continental margin of Laurasia. Longitudinal sedimentary tapering of pre-orogenic sediments correlates strongly with curvature of the submarine frontal accretionary belt in the South China Sea margin, suggesting that pre-orogenic sediment thickness is the major control on the geometry of frontal structures. The preexisting South China Sea slope that lies obliquely in front of the advancing accretionary wedge has impeded the advancing of frontal folds resulting in a successive termination of folds against and along strike of the South China Sea slope. The existence of the South China Sea slope also leads the strike of impinging folds with NNW-trend to turn more sharply to a NE-strike, parallel to strike of the South China Sea slope. Analysis shows that the pre-orogenic mechanical/crustal heterogeneities and seafloor morphology exert strong controls on the thrust-belt development in the incipient Taiwan arc-continent collision zone. In accretionary wedges, seismicity activating superimposed thrusts may drive methane and oil upraising from the upper crust. Mechanical models that treat accretionary complexes as critically tapered wedges of sediment demonstrate that pore pressure controls their taper angle by modifying basal and internal shear strength. Results from some studies show that pore pressure in accretionary wedges can be viewed as a dynamically maintained response to factors which drive pore pressure (source terms) and those that limit flow (permeability and drainage path length). Sediment permeability and incoming sediment thickness are the most important factors, whereas fault permeability and the partitioning of sediment have a small effect. In one such study, it was found that as sediment permeability is increased, pore pressure decreases from near-lithostatic to hydrostatic values and allows stable taper angles to increase from ~2.5° to 8°–12.5°. With increased sediment thickness (from ), increased pore pressure drives a decrease in stable taper angle from 8.4°–12.5° to <2.5–5°. In general, low-permeability and thick incoming sediment sustain high pore pressures consistent with shallowly tapered geometry, whereas high-permeability and thin incoming sediment should result in steep geometry. Active margins characterized by a significant proportion of fine-grained sediment within the incoming section, such as northern Antilles and eastern Nankai, exhibit thin taper angles, whereas those characterized by a higher proportion of sandy turbidites, such as Cascadia, Chile, and Mexico, have steep taper angles. Observations from active margins also indicate a strong trend of decreasing taper angle (from >15° to <4°) with increased sediment thickness (from <1 to 7 km). Rapid tectonic loading of wet sediment in accretionary wedges is likely to cause the fluid pressure to rise until it is sufficient to cause dilatant fracturing. Dewatering of sediment that has been underthrust and accreted beneath the wedge can produce a large steady supply of such highly overpressured fluid. Dilatant fracturing will create escape routes, so the fluid pressure is likely to be buffered at the value required for the transition between shear and oblique tensile (dilatant) fracture, which is slightly in excess of the load pressure if the maximum compression is nearly horizontal. This in turn buffers the strength of the wedge at the cohesive strength, which is not pressure-dependent, and will not vary greatly throughout the wedge. Near the wedge front the strength is likely to be that of the cohesion on existing thrust faults in the wedge. The shear resistance on the base of the wedge will also be fairly constant and related to the cohesive strength of the weak sediment layer that acts as the basal detachment. These assumptions allow the application of a simple plastic continuum model, which successfully predicts the observed gently convex taper of accretionary wedges. Pelayo and Weins have postulated that some tsunami events have resulted from rupture through the sedimentary rock along the basal decollement of an accretionary wedge. Backthrusting of the rear of the accretionary wedge, arcward over the rocks of the forearc basin, is a common aspect of accretionary tectonics. An older assumption that backstops of accretionary wedges dip back toward the arc, and that accreted material is emplaced below such backstops, is contradicted by observations from many active forearcs that indicate (1) backthrusting is common, (2) forearc basins are nearly ubiquitous associates of accretionary wedges, and (3) forearc basement, where imaged, appears to diverge from the sedimentary package, dipping under the wedge while the overlying sediments are often lifted up against it. Backthrusting may be favored where relief is high between the crest of the wedge and the surface of the forearc basin because the relief must be supported by shear stress along the backthrust. Examples Currently active wedges Mediterranean Ridge – part of the active collision zone between the African and Eurasian plates Barbados Ridge – the South American plate is subducting beneath Caribbean plate Nankai accretionary complex – the Philippine Sea plate is subducting beneath the Amurian microplate. In recent years, this is the site of attention for studying the temperature of subseafloor life and underground hot fluids in subducting zones. Exhumed ancient wedges Chilean Coast Range between 38°S and 43°S (Bahía Mansa Metamorphic Complex). Calabrian Accretionary Wedge in the Central Mediterranean – The Neogene tectonics of the central Mediterranean are related to the subduction and trench rollback of the Ionian basin under Eurasia, causing the opening of the Liguro-Provençal and Tyrrhenian back-arc basins and the formation of the Calabrian accretionary wedge. The Calabrian accretionary wedge is a partially submerged accretionary complex located in the Ionian offshore and laterally bounded by the Apulia and Malta escarpments. The Olympic Mountains located in Washington State. The mountains began to form about 35 million years ago when the Juan de Fuca plate collided with and was forced (subducted) under the North American plate. Kodiak Shelf in the Gulf of Alaska – The geology of the Chugach National Forest is dominated by two major lithologic units, the Valdez Group (Late Cretaceous) and the Orca Group (Paleocene and Eocene). The Valdez Group is part of a 2,200-km-long by 100-km-wide belt of Mesozoic accretionary complex rocks called the Chugach terrane. This terrane extends along the Alaska coastal margin from Baranof Island in southeastern Alaska to Sanak Island in southwestern Alaska. The Orca Group is part of an accretionary complex of Paleogene age called the Prince William terrane that extends across Prince William Sound westward through the Kodiak Island area, underlying much of the continental shelf to the west Neogene accretionary wedge off Kenai Peninsula, Alaska – Subduction accretion and repeated terrane collision shaped the Alaskan convergent margin. The Yakutat Terrane is currently colliding with the continental margin below the central Gulf of Alaska. During the Neogene the terrane's western part was subducted after which a sediment wedge accreted along the northeast Aleutian Trench. This wedge incorporates sediment eroded from the continental margin and marine sediments carried into the subduction zone on the Pacific plate. The Franciscan Formation of California – Franciscan rocks in the Bay Area range in age from about 200 million to 80 million years old. The Franciscan Complex is composed of a complex amalgamation of semi-coherent blocks, called tectonostratigraphic terranes, that were episodically scraped from the subducting oceanic plate, thrust eastward, and shingled against the western margin of North America. This process formed a stacking sequence in which the structurally highest rocks (on the east) are the oldest, and in which each major thrust wedge to the west becomes younger. Within each of the terrane blocks, however, the rocks become younger upsection, but the sequence may be repeated multiple times by thrust faults. The Apennines in Italy are largely an accretionary wedge formed as a consequence of subduction. This region is tectonically and geologically complex, involving both subduction of the Adria micro-plate beneath the Apennines from east to west, continental collision between the Eurasia and Africa plates building the Alpine mountain belt further to the north and the opening of the Tyrrhenian basin to the west. Carpathian Flysch Belt in Bohemia, Slovakia, Poland, Ukraine and Romania represent Cretaceous to Neogene thin-skinned zone of Carpathian thrustbelt, which is thrust over the Bohemian Massif and East European Platform. Represents a continuation of Alpine Rhenodanubian Flysch of Penninic Unit.	Physical sciences	Tectonics	Earth science
3272528	https://en.wikipedia.org/wiki/Four-eyed%20fish	Four-eyed fish	The four-eyed fishes are a genus, Anableps, of fishes in the family Anablepidae. They have eyes raised above the top of the head and divided in two different parts, so that they can see below and above the water surface at the same time. The optomotor response or OMR has been used as a test to investigate potential differential visual processing in Anableps on normal versus ‘blinded’ fish (the eyes are actually covered—not physically blinded). It was found that the OMR does exist in Anableps and that the strength of this response is dependent on the visual field being tested—a stronger OMR was seen as a result of visual stimulation from the aerial environment. Like their relatives, the onesided livebearers, four-eyed fishes mate only on one side, right-"handed" males with left-"handed" females and vice versa. These fish inhabit fresh and brackish water and are only rarely coastal marine. They originate in lowlands in southern Mexico to Honduras and northern South America, but can also be found on the island of Trinidad located to the east of Venezuela. Species There are currently three recognized species in this genus: Anableps anableps (Linnaeus, 1758) (Largescale foureyes) Anableps dowei T. N. Gill, 1861 (Pacific foureyed fish) Anableps microlepis J. P. Müller & Troschel, 1844 (Foureyes) Physical characteristics The maximum length of four-eyed fishes is up to 32 cm TL in A. microlepis, making this species the largest in the order Cyprinodontiformes. Four-eyed fish have only two eyes, but the eyes are specially adapted for their surface-dwelling lifestyle. In early development, the four-eyed fish’s frontal bone expands dorsally allowing the eyes to be positioned on top of their head and appear bulging. This allows the fish to simultaneously see above and below the water as it floats at the surface. The eyes are divided into dorsal and ventral halves, separated by a pigmented strip of tissue. Each eye has two pupils and two corneas filtering light onto one lens, refracting onto separate hemiretinas and processed through one optic disc. The upper (dorsal) half of the eye is adapted for vision in air, the lower (ventral) half for vision in water. The lens of the eye also changes in thickness top to bottom to compensate for the difference in the refractive indices of air versus water. The ventral hemiretina is characterized by thicker cell layers containing more sensory neurons and an increased visual acuity compared to the dorsal hemiretina. Four-eyed fish are livebearers. Along with their sister genus Jenynsia they mate on one side only, right-"handed" males with left-"handed" females and vice versa. The male has specialized anal rays which are greatly elongated and fused into a tube called a gonopodium associated with the sperm duct which he uses as an intromittent organ to deliver sperm to the female. Behavior Four-eyed fish spend most of their time at the surface of the water. Their diet mostly consists of terrestrial insects which are readily available at the surface, however they may consume other foods such as other invertebrates, diatoms, and small fishes. The fish will group differently depending on the species. A. anableps commonly congregates in schools. A. microlepis also is gregarious, but restricts its schools to about a dozen individuals; it is also recorded to be found alone or as couples. A. anableps is also known for the ability to survive out of water when exposed to air, especially during low tide.	Biology and health sciences	Acanthomorpha	Animals
3276702	https://en.wikipedia.org/wiki/Bicycle%20boulevard	Bicycle boulevard	A bicycle boulevard, sometimes referred to as a neighborhood greenway, neighborway, neighborhood bikeway or neighborhood byway is a type of bikeway composed of a low-speed street which has been "optimized" for bicycle traffic. Bicycle boulevards discourage cut-through motor-vehicle traffic but may allow local motor-vehicle traffic at low speeds. They are designed to give priority to bicyclists as through-going traffic. They are intended as a low-cost, politically popular way to create a connected network of streets with good bicyclist comfort and/or safety. Bicycle boulevards attempt to achieve several goals: discouragement of non-local motor vehicle traffic; low speed limits; low motor-vehicle traffic volumes; free-flow travel for bikes by assigning the right-of-way to the bicycle boulevard at intersections wherever possible; traffic control to help bicycles cross major arterial roads; a distinctive look and/or ambiance such that cyclists become aware of the existence of the bike boulevard and motorists are alerted that the street is a priority route for bicyclists; and, enhanced environment due to the promotion of bicycle usage. These bikeway design elements are intended to appeal to casual, risk-averse, inexperienced and younger bicyclists who would not otherwise be willing to cycle with motor vehicle traffic. Compared to a bike path or rail trail, a bicycle boulevard is also a relatively low-cost approach to appealing to a broader cycling demographic. Features A bicycle boulevard is generally marked with a sign at the beginning and the end of the bicycle boulevard. Also necessary for the road to be called a bicycle boulevard is coloring; in the Netherlands, the parts of the road where the cyclists ride on is marked in red (same color as used for segregated cycle facilities in the Netherlands). These sections of the road are called . Motorists also ride on this section, yet also have a non-colored part of the road which they can drive on with one half (two wheels) of the car when they wish to pass a cyclist. Bicycle boulevards may use a variety of traffic calming elements to achieve a safe environment. This makes it difficult for motorists to use the street at a high speed. However, they do not block access to motor vehicles completely (i.e. using bollards) which would designate the route as segregated cycle facilities rather than a bicycle boulevard. Some bicycle boulevards have higher road surface standards than other residential streets, and encourage riders to use the full lane, encouraging parity between bicycles and motor vehicles. Discouraging non-local motor vehicle traffic Permeable barriers such as bollards are sometimes used to allow cycling traffic to continue through while diverting motorized traffic from using the street as a through street. Locations Road designs of bicycle boulevards can be found in the United States, Canada (Vancouver, Saskatoon, Winnipeg), the Netherlands, Germany, Belgium, Denmark, France, Spain and New Zealand United States Bicycle boulevards can be found in a growing number of United States cities, including: Arizona: Tucson California: Palo Alto, Berkeley, Emeryville, San Jose, San Luis Obispo, Long Beach Florida: Gainesville Kansas: Manhattan Michigan: Kalamazoo Minnesota: Minneapolis, Saint Paul Missouri: Columbia New Mexico: Albuquerque North Carolina: Wilmington Oregon: Portland, Eugene, and Bend Oklahoma: Tulsa Washington: Seattle Wisconsin: Madison Palo Alto established the first bicycle boulevard in the United States. It was named for Ellen Fletcher, a Holocaust survivor and one of America's first bike activists. It was an overall success minus a few complaints from local residents that were clarified through city council meetings and test pilots. In Berkeley, boulevards are mostly residential streets, but some sections pass through commercial areas. Generally, there are few cars on these streets, in large part because of the pre-existing traffic calming devices that slow and/or divert traffic. Bicycle boulevards may or may not have bicycle lanes. In Minneapolis, a grant from the federal government within the Non-Motorized Pilot Program helped to build a bike boulevard on Bryant Avenue and the planning of others. Similarly in Columbia, the Non-Motorized Pilot Program project helped fund the first bike boulevard in Missouri along Ash and Windsor Streets. At least one other was planned. In Wilmington, help from a Fit Community 2009 grant through the North Carolina Health and Wellness Trust Fund enabled the City of Wilmington to construct North Carolina's first bicycle boulevard. The Ann Street Bicycle Boulevard runs from South Water Street to South 15th Street and serves as part of the much longer River to the Sea Bikeway, which connects downtown Wilmington to Wrightsville Beach. In Portland, a $600 million 20-year plan (2010–2030) has the goal of making 25 percent of trips in the city be by bicycle through the establishment of of new bikeways; one of the projects within the plan is to combine the work on street features that reduce stormwater runoff with the construction of curb extensions and other components of bicycle boulevards. In Albuquerque, a city with more than of on-street bicycle facilities and multi-use trails, the grand opening of the first bicycle boulevard in New Mexico was held on April 14, 2009. The bicycle boulevard runs from San Mateo Blvd SE, west along Silver Ave SE/SW to 14th St SW. It then continues north on 14th St to Mountain Rd NW. The last leg continues west on Mountain Rd NW to the Paseo del Bosque Recreation Trail which parallels the Rio Grande. In Madison, the first full bicycle boulevard spans East Mifflin Street in Madisons Tenney-Lapham Neighborhood, a second spans the entire length of Kendall Avenue in University Heights and the Regent Neighborhood. In Seattle, the city is implementing a city-wide network of "Neighborhood Greenways". The work is being carried out with the aid and cooperation of the non-profit "Seattle Neighborhood Greenways". US naming conventions The City of Berkeley, California, is credited with coining the phrase bicycle boulevard in the late 1980s, but not every jurisdiction has adopted this term. In November 2011, the City of Boston began to use the term neighborways instead of bicycle boulevards. This added to a growing list of terms for bicycle boulevards since Portland has been calling them neighborhood greenways; Seattle has followed the same convention. Other terms for bicycle boulevards in the US include: Cyclestreets Bike boulevards Quiet streets Neighborhood byways Bicycle-friendly streets Bicycle-friendly corridors Bicycle parkways Neighborhood parkways Bicycle greenways Netherlands In the Netherlands, ('cycle streets') have a similar road design — although most residential streets in the Netherlands which do not have on-road bike lanes or segregated bike lanes would fit the American definition of bicycle boulevards. A can link dedicated bike-only paths, service roads, and other types of bike-friendly street configurations to complete a route. (Extensive information has been compiled about these facilities at the Pedal Portland blog and the Northeastern University webpage.) In Amsterdam for example, by 2005 about 40% of journeys were by bicycle and transport planners at the (Infrastructure Traffic and Transport Directorate) have adopted a bicycle policy that blends many different bike-friendly street designs such as segregated bicycle lanes, on-road bicycle lanes and fietsstraten. The general concept is that cyclists can integrate relatively safely with vehicular traffic that is travelling at, or below, but that segregated bike lanes should be installed along roads with a higher speed limit. With these, and many other, bike-friendly policies in place, Amsterdam has the highest rate of cycling of any capital city in the world. Cycle streets are also on the rise in other cities within the country, including Utrecht. Germany In Germany a comparable road design is called ('bicycle road'), introduced into the Highway Code in 1997. Any other vehicles are prohibited unless marked with an additional sign. Belgium In Belgium, the (in Dutch/Flemish) or (in French), was introduced into the Highway Code with effect from 13 February 2012. One had earlier been introduced in the Visserij in Ghent (Gent) in the summer of 2011. The first one in Brussels appeared in 2013 on a service road alongside Avenue Louise. The OpenStreetMap wiki and also the several locations on this subject may be of interest to reader. Denmark In Denmark, the first ('cycle street') was opened in 2011 in Aarhus. Since then cycle streets have been implemented in several cities across the country. Almost all Danish cycle streets allow motorized vehicles to drive on them, although some might be one way only. The speed limit is 50km/h although the law states that drivers should limit their speed to that of cyclists, normally under 30km/h France In France, the equivalent road design is called ('bike street') or ('cycle street'). The cities of Strasbourg (2017), Bordeaux (2018), Dijon (2019), Rennes (2023), Paris (2023), Lille (2023) and Lyon (2024) are among the first to test it out. Spain In Spain, cycle streets are known as . Sweden In Sweden cycle streets are known as Cykelgata and was introduced in december 2020. New Zealand In New Zealand, bicycle boulevards are generally designated as 'neighbourhood greenways', although Auckland refers to them as local paths to avoid confusion with its off-road greenways network. Christchurch was the first city to implement a number of neighbourhood greenway sections as part of its Major Cycle Routes programme, including the Rapanui–Shag Rock Stage 1 through Linwood.	Technology	Road infrastructure	null
17992940	https://en.wikipedia.org/wiki/Sodium%20chromate	Sodium chromate	Sodium chromate is the inorganic compound with the formula Na2CrO4. It exists as a yellow hygroscopic solid, which can form tetra-, hexa-, and decahydrates. It is an intermediate in the extraction of chromium from its ores. Production and reactivity It is obtained on a vast scale by roasting chromium ores in air in the presence of sodium carbonate: 2Cr2O3 + 4 Na2CO3 + 3 O2 → 4 Na2CrO4 + 4 CO2 This process converts the chromium into a water-extractable form, leaving behind iron oxides. Typically calcium carbonate is included in the mixture to improve oxygen access and to keep silicon and aluminium impurities in an insoluble form. The process temperature is typically around 1100 °C. For lab and small scale preparations a mixture of chromite ore, sodium hydroxide and sodium nitrate reacting at lower temperatures may be used (even 350 C in the corresponding potassium chromate system). Subsequent to its formation, the chromate salt is converted to sodium dichromate, the precursor to most chromium compounds and materials. The industrial route to chromium(III) oxide involves reduction of sodium chromate with sulfur. Acid-base behavior It converts to sodium dichromate when treated with acids: 2 Na2CrO4 + 2HCl → Na2Cr2O7 + 2NaCl + H2O Further acidification affords chromium trioxide: Na2CrO4 + H2SO4 → CrO3 + Na2SO4 + H2O Uses Aside from its central role in the production of chromium from its ores, sodium chromate is used as a corrosion inhibitor in the petroleum industry. It is also a dyeing auxiliary in the textile industry. It is a diagnostic pharmaceutical in determining red blood cell volume. In organic chemistry, sodium chromate is used as an oxidant, converting primary alcohols to carboxylic acids and secondary alcohols to ketones. Sodium chromate is a strong oxidizer. Safety As with other Cr(VI) compounds, sodium chromate is carcinogenic. The compound is also corrosive and exposure may produce severe eye damage or blindness. Human exposure further encompasses impaired fertility, heritable genetic damage and harm to unborn children.	Physical sciences	Metallic oxyanions	Chemistry
17999576	https://en.wikipedia.org/wiki/Palaeomerycidae	Palaeomerycidae	The Palaeomerycidae is an extinct family of Neogene ruminants belonging to the infraorder Pecora. Palaeomerycids lived in Europe and Asia exclusively during the Miocene, coevolving with cervids, bovids, moschids, and tragulids there as part of a dramatic radiation of ruminants by the early Miocene. Dromomerycids are sometimes considered to be subfamilies of the Palaeomerycidae, but recent research brought doubt to this, arguing that the dromomerycids lack the sutures on the skull roof that giraffomorphs (Giraffidae, Palaeomerycidae, Climacoceratidae) have for ossicone features. The similar resemblances of the two families could be the result of parallel evolution. Description Palaeomerycids were a group of horned, long-legged and massive ruminants that could attain a weight of . One of the first known members of this group, Palaeomeryx, was thought to be a hornless form distantly related to the Giraffidae before paleontologist Miguel Crusafont found remains of Triceromeryx in middle Miocene Spain. This Palaeomeryx-like form carried two ossicones over its orbits that were straight and short, similar to those of true giraffids. However, the most striking feature of Triceromeryx was the third, Y-shaped appendage that prolonged the occipital bone at the back of the skull. Discoveries during the 1980s and 1990s showed a surprising variety in these occipital appendages. Ampelomeryx, a genus of palaeomerycids found at the early Miocene sites of Els Casots, Valles-Penedes Basin, Spain, and Montréal-du-Gers, Gers, France, had a three-horned system of appendages similar to those of Triceromeryx. These appendages were, however, quite different, with the paired appendages extending laterally over the orbits, flat and wide, forming an eye-shade, while the third spectacular posterior appendage was about long. Another species of Triceromeryx, T. conquensis found in La Retama in Spain, showed an even more spectacular appendage — instead of a Y-shaped structure, its posterior appendage was T-shaped with the lateral branches expanding toward the front. In primitive members of the group (e.g. Ampelomeryx), this appendage was a posterior expansion of the occipital bone lying close to the powerful muscles supporting the skull in a normal position, thus suggesting that this appendage was actually used for fighting between males during the breeding season. The reduced shapes of the flat and laterally oriented appendages of later species suggests that these were not used in active fighting, instead forming a function of passive display. As a group, the palaeomerycids appear to have formed a successful part of an independent radiation of horned ruminants that diversified into a variety of forms during the early to middle Miocene, with a geographic range reaching from Spain to China. Taxonomy The Palaeomerycidae were named by Lydekker (1883). The type genus is Palaeomeryx. The family was assigned to the Artiodactyla by Hulbert and Whitmore (2006), and to Cervoidea by Carroll (1988), Sach and Heizmann (2001) and Prothero and Liter (2007). Classification Palaeomerycinae Ampelomeryx Palaeomeryx Triceromeryx Xenokeryx Tauromeryx	Biology and health sciences	Other artiodactyla	Animals
1712680	https://en.wikipedia.org/wiki/Percent%20sign	Percent sign	The percent sign (sometimes per cent sign in British English) is the symbol used to indicate a percentage, a number or ratio as a fraction of 100. Related signs include the permille (per thousand) sign and the permyriad (per ten thousand) sign (also known as a basis point), which indicate that a number is divided by one thousand or ten thousand, respectively. Higher proportions use parts-per notation. Correct style Form and spacing English style guides prescribe writing the percent sign following the number without any space between (e.g. 50%). However, the International System of Units and ISO 31-0 standard prescribe a space between the number and percent sign, in line with the general practice of using a non-breaking space between a numerical value and its corresponding unit of measurement. Other languages have other rules for spacing in front of the percent sign: In Czech and in Slovak, the percent sign is spaced with a non-breaking space if the number is used as a noun. In Czech, no space is inserted if the number is used as an adjective (e.g. "a 50% increase"), whereas Slovak uses a non-breaking space in this case as well. In Croatian, the percent sign is spaced with a non-breaking space. In Finnish, the percent sign is always spaced, and a case suffix can be attached to it using the colon (e.g. 50 %:n kasvu 'an increase of 50%'). In French, the percent sign must be spaced with a non-breaking space. According to the Real Academia Española, in Spanish, the percent sign should be spaced now, despite the fact that it is not the linguistic norm. Despite that, in North American Spanish (Mexico and the US), several style guides and institutions either recommend the percent sign be written following the number without any space between or do so in their own publications in accordance with common usage in that region. In Russian, the percent sign is rarely spaced, contrary to the guidelines of the GOST 8.417-2002 state standard. In Chinese, the percent sign is almost never spaced, probably because Chinese does not use spaces to separate characters or words at all. According to the Swedish Language Council, the percent sign should be preceded by a space in Swedish, as all other units. In German, the space is prescribed by the regulatory body in the national standard DIN 5008. In Turkish and some other Turkic languages, the percent sign precedes rather than follows the number, without an intervening space. In Persian texts, the percent sign may either precede or follow the number, in either case without a space. In Arabic, the percent sign follows the number; as Arabic is written from right to left, this means that the percent sign is to the left of the number, usually without a space. In Hebrew, the percent sign is written to the right of the number, just as in English, without an intervening space. This is because numbers in Hebrew (which otherwise is written from right to left) are written from left to right, as in English. In Dutch, the official rule (NBN Z 01-002) is to place a space between the number and the sign (e.g. "een stijging van 50 %"), but most of the time, the space is missing (e.g. "een stijging van 50%"). Usage in text It is often recommended that the percent sign only be used in tables and other places with space restrictions. In running text, it should be spelled out as percent or per cent (often in newspapers). For example, not "Sales increased by 24% over 2006" but "Sales increased by 24 percent over 2006". Evolution Prior to 1425, there is no known evidence of a special symbol being used for percentages. The Italian term per cento, "for a hundred", was used as well as several different abbreviations (e.g. "per 100", "p 100", "p cento"). Examples of this can be seen in the 1339 arithmetic text (author unknown) depicted below. The letter p with its descender crossed by a horizontal or diagonal stroke (to indicate abbreviation) conventionally stood for per, por, par, or pur in Medieval and Renaissance palaeography. At some point, a scribe used the abbreviation pc with a tiny loop or circle (depicting the ending -o used in Italian ordinals, as in , ; it is analogous to the English -th as in 25th). This appears in some additional pages of a 1425 text which were probably added around 1435. The pc with a loop eventually evolved into a horizontal fraction sign by 1650 (see below for an example in a 1684 text) and thereafter lost the per. In 1925, D. E. Smith wrote, "The solidus form () is modern." Encodings Unicode The Unicode code points are: (HTML %, &percnt;), , which has the circles replaced by square dots set on edge, the shape of the digit 0 in Eastern Arabic numerals. , (, from English "percent") in one square character. ASCII The ASCII code for the percent character is 37, or 0x25 in hexadecimal. Other uses In computers Names for the percent sign include percent sign (in ITU-T), mod, grapes (in hacker jargon), and the humorous double-oh-seven (in INTERCAL). In computing, the percent character is also used for the modulo operation in programming languages that derive their syntax from the C programming language, which in turn acquired this usage from the earlier B. In the textual representation of URIs, a % immediately followed by a 2-digit hexadecimal number denotes an octet specifying (part of) a character that might otherwise not be allowed in URIs (see percent-encoding). In SQL, the percent sign is a wildcard character in "LIKE" expressions, for example will fetch all records whose names start with "". In TeX (and therefore also in LaTeX) and PostScript, and in GNU Octave and MATLAB, a % denotes a line comment. In BASIC, Visual Basic, ASP, and VBA a trailing % after a variable name marks it as an integer. In ASP, the percent sign can be used to indicate the start and end of the ASP code In Perl % is the sigil for hashes. In many programming languages' string formatting operations (performed by functions such as printf and scanf), the percent sign denotes parts of the template string that will be replaced with arguments. (See printf format string.) In Python and Ruby the percent sign is also used as the string formatting operator. In the command processors COMMAND.COM (DOS) and CMD.EXE (OS/2 and Windows), ,... stand for the first, second,... parameters of a batch file. %0 stands for the specification of the batch file itself as typed on the command line. The % sign is also used similarly in the FOR command. %VAR1% represents the value of an environment variable named VAR1. Thus: sets a new value for PATH, that being the old value preceded by "c:\;". Because these uses give the percent sign special meaning, the sequence %% (two percent signs) is used to represent a literal percent sign, so that: would set PATH to the literal value "c:\;%PATH%". In the C Shell, % is part of the default command prompt. In linguistics In linguistics, the percent sign is prepended to an example sentence or other string to show that it is judged well-formed (grammatical) by some speakers and ill-formed by others. This may be due to differences in dialect or even individual idiolects. This use is similar to those of the asterisk to mark ill-formed strings, the question mark to mark strings where well-formedness is unclear, and the number sign to mark strings that are syntactically well-formed but semantically or pragmatically nonsensical.	Mathematics	Basics	null
1714377	https://en.wikipedia.org/wiki/Liopleurodon	Liopleurodon	Liopleurodon (; meaning 'smooth-sided teeth') is an extinct genus of carnivorous pliosaurid pliosaurs that lived from the Callovian stage of the Middle Jurassic to the Kimmeridgian stage of the Late Jurassic period (c. 166 to 155 mya). The type species is L. ferox, which is probably the only valid species. Some studies also include the second species L. pachydeirus, but this latter is considered as a probable junior synonym of L. ferox due to its lack of viable diagnosis. As the holotype specimen of L. ferox consists of a single tooth preserving questionable distinctive features, recent studies therefore recommend the necessary identification of a neotype in order to preserve the validity of the genus. Numerous fossil specimens attributed to Liopleurodon, even including numerous skeletons, have been discovered in Europe, Russia, and Mexico. Other additional species were even proposed, but these are currently seen as coming from other pliosaurid genera. Liopleurodon is a representative of the Thalassophonea clade, a derived group of pliosaurids characterized by a short neck and a large elongated skull. In 1999, the size of Liopleurodon was greatly exaggerated in the BBC documentary series Walking with Dinosaurs, depicted as reaching in length. However, the different attributed specimens show that the animal could reach a size ranging from long, with some researchers estimating a maximum length of approximately . Various studies suggest that Liopleurodon would have been an ambush predator, feeding on fish, cephalopods and other marine reptiles. Research history Even before Liopleurodon was named, material likely belonging to it was described. In 1841, Hermann von Meyer named the species Thaumatosaurus oolithicus based on a fragmentary specimen consisting of partial teeth, skull elements, vertebrae, and ribs from deposits in Württemberg, Germany, possibly dating to the Oxfordian. However, this material is nondiagnostic, lacking distinguishing features. Johann Andreas Wagner published a description of a large plesiosaur tooth from Bavaria, Germany, in 1852, assigning it to a new species that he named Pliosaurus giganteus. However, in 1824, William Conybeare had named a species of Plesiosaurus, Plesiosaurus giganteus, and this species was later viewed as a synonym of either Pliosaurus brachydeirus or P. brachyspondylus by following authors. Since the name Pliosaurus giganteus had been used prior to Wagner's publication, Wagner's name is invalid due to preoccupation. In 1838, Hermann von Meyer applied the name Ischyrodon meriani to a large tooth from Oxfordian-aged rocks in Fricktal, Switzerland. This tooth lacks identifying characteristics, and therefore it is not clear what it belonged to, although Lambert Beverly Tarlo noted the possibility of it pertaining to Liopleurodon in 1960. A 2022 study by Daniel Madzia and colleagues noted that while the tooth likely came from Liopleurodon or something similar, there was too little information available to make a confident assignment, so they treated Ischyrodon as a nomen dubium. In 1860, Hermann Trautschold assigned the name Pliosaurus giganteus to a small tooth now thought to pertain to Liopleurodon. However, as the name Pliosaurus giganteus had already been used twice by this point, Trautschold's name is also invalid. The genus name Liopleurodon was coined by Henri Émile Sauvage in 1873. Sauvage named three species which he assigned to this genus, each based on a single tooth. One tooth, its crown measuring long, was found near Boulogne-sur-Mer, France, in layers dating from the Callovian, and was named Liopleurodon ferox. Another from Charly, France, measuring long and with a crown length of , was named Liopleurodon grossouvrei. The third, discovered near Caen, France, was originally attributed to Poikilopleuron bucklandi by Eudes Deslongchamps. While the tooth could have come from the megalosaur, Sauvage considered this identity unsubstantiated, and assigned it to the species Liopleurodon bucklandi. Sauvage did not ascribe the genus to any particular group of reptiles in his descriptions. The genus name Liopleurodon derives from Ancient Greek , "smooth"; , "side" or "rib"; and , "tooth", all meaning "smooth-sided tooth", in reference to the very contrasting dentition of the animal. The specific name derives from Latin ferox, meaning "fierce", in reference to the large size of the teeth. However, in 1880, Sauvage synonymized Liopleurodon with Polyptychodon, noting that it was similar to this genus, but distinct from Plesiosaurus and Pliosaurus. In 1888, Richard Lydekker, after studying some teeth attributable to Liopleurodon ferox in the Leeds Collection, concluded that they were so similar to those of Pliosaurus that they should be placed in that genus. These teeth had been collected by Alfred Leeds from the Oxford Clay Formation, near Peterborough, England. In 1869, Harry Govier Seeley had applied the name Pliosaurus pachydeirus to a series of cervical (neck) vertebrae representing the first 17 in the neck from the Oxford Clay Formation near Great Gransden. Other than its large size, Seeley provided no distinguishing characteristics. The specific name pachydeirus means "stout neck", due to its vertebral morphology. Lydekker stated in 1888 that the neck described by Seeley probably belonged to Pliosaurus ferox. W. Kiprijanoff named Thaumatosaurus mosquensis in 1883 based on remains including teeth, vertebrae, and limb bones from Oxfordian-aged rocks in the Moscow Basin of Russia. However, in 1889, Lydekker considered this species to be a probable junior synonym of P. ferox. In 1905, John Frederick Blake described two teeth from Rushden, England, similar to those of other Liopleurodon ferox specimens, though from older strata than those from Peterborough. He noted that the teeth were quite different from those of Pliosaurus, while the bones were dissimilar to those of Polyptychodon. Since the species couldn't be assigned to either genus, he recommended reinstating the name Liopleurodon. After considering Liopleurodon to be a subgenus of Pliosaurus, N. Bogolubov also listed the two genera as distinct in 1912. When Lydekker had first visited the collection of Alfred Leeds (known as the Leeds Collection), the only remains of Liopleurodon in his collection were teeth. However, since then, Alfred Leeds, as well as his brother Charles Edward Leeds, had collected many more specimens of Liopleurodon, including skulls and much of the postcranial skeleton. Charles William Andrews described the anatomy of the marine reptile specimens of the Leeds Collection acquired by the British Museum of Natural History in two volumes, the first published in 1910 and the second in 1913. He described the Liopleurodon specimens in the second volume, though considered them to belong to Pliosaurus. Hermann Linder also described specimens of Liopleurodon ferox in 1913. One of these was a poorly preserved partial skeleton excavated from the Oxford Clay of Fletton, England, housed in Institut für Geowissenschaften, University of Tübingen. The skeleton was mounted and missing regions were restored with material from other Liopleurodon specimens. Like Andrews, Linder also considered L. ferox to be a species of Pliosaurus. Additionally, Linder described some skulls from Fletton housed at both the University of Tübingen and the State Museum of Natural History Stuttgart as specimens of P. grandis. Linder also assigned a nearly complete paddle to Pliosaurus sp. All of these specimens have since been assigned to Liopleurodon with varying degrees of confidence, though the skull Linder attributed to P. grandis that was housed in Stuttgart was destroyed during World War II. In 1934, Friedrich von Huene described a partial skeleton from Swabia, Germany. He also used Pliosaurus ferox instead of Liopleurodon ferox. In 1939, Alexandre Bigot used Pliosaurus ferox as well, assigning some teeth from France to this species. Lambert Beverly Halstead, then known as Tarlo, published a review of Upper Jurassic pliosaurid taxonomy in 1960. He considered Liopleurodon to be distinct from Pliosaurus, noting major differences between the mandibles of the two genera. In addition to the type species L. ferox, Tarlo also considered Pliosaurus pachydeirus to be a valid species within Liopleurodon, L. pachydeirus, noting that the two species had differences in their teeth and cervical vertebrae. L. grossouvrei was not considered valid, though it was tentatively retained for teeth from the Kellaways Formation. In 1971, Halstead published another paper about Jurassic pliosaurids, this time focusing on Pliosaurus rossicus, a species he was formerly unwilling to consider valid, due to a lack of information. After reviewing its anatomy, he considered it valid, though assigned it to Liopleurodon instead, based on its short mandibular symphysis. Halstead also considered Pliosaurus macromerus, which he had previously considered to belong to its own genus, Stretosaurus, to instead be a species of Liopleurodon, despite its irregularly-shaped scapula (although this was later discovered to be an ilium). In 1992, Martill identified a fragmentary specimen belonging to a young individual, PETCM R.296, as cf. Liopleurodon sp.; the specimen was found to have at least 7 gastroliths in its stomach and soft tissues, although the specific features of the latter cannot be observed due to poor preservation. In a 2001 dissertation, Leslie F. Noè argued that L. pachydeirus was not diagnostic, and that L. ferox was the only valid species of Liopleurodon. The teeth of mounted skeleton in Tübingen, which Tarlo had attributed to L. pachydeirus, showed distinctive characteristics of L. ferox, indicating that cervical vertebrae are not useful for differentiating species, as argued by David S. Brown in 1981. While Tarlo had considered differences in tooth morphology to be diagnostic, Noè instead considered it to be individual variation. Noè also removed L. macromerus and L. rossicus from the genus, citing differences in tooth shape and mandibular symphysis length. The former species was tentatively placed back in Pliosaurus, while the latter was thought to warrant a new genus. Liopleurodon fossils have been found mainly in England and France. Fossil specimens that are contemporary (Callovian-Kimmeridgian) with those from England and France referrable to Liopleurodon are known from Germany. In 2013, Roger Benson and colleagues considered both "L." macromerus and "L." rossicus to belong to Pliosaurus. They also considered Liopleurodon to be restricted to the Middle Jurassic. In 2015, Jair Israel Barrientos-Lara and colleagues described two pliosaurid fossils found near the town of Tlaxiaco in Oaxaca, Mexico. These fossils were extracted from Kimmeridgian deposits in the Sabinal Formation, and one of them, the partial front end of a snout, was attributable to Liopleurodon, though the researchers considered the remains too fragmentary to provide a species-level identification. Liopleurodon grossouvrei, although synonymized with Pliosaurus andrewsi by most authors, was considered to potentially be a distinct genus in its own right by Davide Foffa and colleagues in 2018, given its differences from P. andrewsi and Liopleurodon ferox. Madzia and colleagues in 2022 noted that the fact that Liopleurodon was named based on a single tooth of dubious distinctiveness is problematic, and that a more complete neotype may need to be designated to preserve the stability of L. ferox. They also stated that further study of the taxon was needed to confirm that the supposed differences between L. ferox and L. pachydeirus were indeed due to individual variation. In 2024, Peggy Vincent and coauthors described a postcranial skeleton measuring around long as preserved from Saint-Laon, France. Found in 1979, this skeleton, informally termed as the Thouarsais specimen, is among the most complete Liopleurodon postcranial skeletons known and the most complete French specimen. Description Liopleurodon ferox first came to the public attention in 1999 when it was featured in an episode of the BBC television series Walking with Dinosaurs, which depicted it as an enormous long and predator; this was based on very fragmentary remains, and considered to be an exaggeration for Liopleurodon, with the calculations of specimens generally considered dubious. Estimating the size of pliosaurs is difficult because not much is known of their postcranial anatomy. The palaeontologist L. B. Tarlo suggested that the pliosaurs’ total body length can be estimated from the length of their skull which he claimed was typically one-seventh of the former measurement. Additional Kronosaurus specimens and a skeleton of L. ferox, GPIT 1754/2, show that the pliosaurs’ skulls were actually about one-fifth of their total body length. One large skull specimen of L. ferox, CAMSMJ.27424, has an estimated total body length of . McHenry estimated that smaller individuals measuring about long would have weighed around based on the specimen NHM R2680. In 2024, Ruizhe Jackevan Zhao estimates that the largest known specimen of Liopleurodon, NHMUK PV R3536, would have reached a length of approximately with a body mass of . Some researchers propose larger estimates of over . Tarlo applied the aforementioned one-seventh ratio of skull length to body length, estimating that the largest known specimen of L. ferox was a little over , though a more typical size range would be from . In the 2023 book Ocean Life in the Time of Dinosaurs, Bardet and colleagues also claimed that some individuals could reach lengths of over . Classification Liopleurodon belongs to clade Thalassophonea, a short necked clade within the Pliosauridae, a family of plesiosaurs, thalassophoneans ranged from the Middle Jurassic to early Late Cretaceous, and have been found worldwide. Liopleurodon was one of the basal taxa from the Middle Jurassic. Differences between these taxa and their relatives from the Upper Jurassic include alveoli count, smaller skull and smaller body size. The following cladogram follows Ketchum and Benson, 2022: Palaeobiology Four strong paddle-like limbs suggest that Liopleurodon was a powerful swimmer. Its four-flipper mode of propulsion is characteristic of all plesiosaurs. A study involving a swimming robot has demonstrated that although this form of propulsion is not especially efficient, it provides very good acceleration—a desirable trait in an ambush predator. Studies of the skull have shown that it could probably scan the water with its nostrils to ascertain the source of certain smells. A fragmentary specimen possibly belonging to a young individual, PETCM R.296, contained numerous hooklets of teuthoid cephalopods, fish bones and a single reptilian tooth in its stomach. Although its exact dietary preference cannot be determined, Martill proposed three suggestions. One possibility is that Liopleurodon could have fed on food supplies that are abundant (i.e. squids), but considering that plesiosaurs and ichthyosaurs were also abundant and that the plesiosaurs' swimming speed is likely very slow compared to squids, this interpretation may be unlikely unless Liopleurodon was an ambush predator. Another possibility is that Liopleurodon may have been an opportunistic feeder, with cephalopod hooklets being representative of the acid resistant residue of its varied diet—skeletal components of various vertebrates that lost to the acid environment of the gut; however, since the thin sections through the gut don't reveal the presence of otoliths (calcium carbonate structure of vertebrates located in the vestibular labyrinth) which are known to occur in the gut of cetaceans, fish may not have been an important part of its diet. The other possibility is that the pliosaur fed on large cephalopod-feeders, with the hooklets representing the residues of the stomach contents of the pliosaur's prey, but there is no firm evidence to this claim. It is also notable that this specimen preserved at least 7 gastroliths, which probably weren't used for grinding based on the well-preserved conditions of the hooklets. It is possible either that the pliosaur accidentally swallowed the stones and they remained in its gut, or that the stones represent the "acid resistant residue from carbonate cemented sandstone."	Biology and health sciences	Prehistoric marine reptiles	Animals
1714999	https://en.wikipedia.org/wiki/Wrestling%20halfbeak	Wrestling halfbeak	The wrestling halfbeak (Dermogenys pusilla) also known as Malayan halfbeak is a species of viviparous halfbeak native to the fresh and brackish waters of rivers and coastal regions in South-East Asia, in Singapore, Indonesia, Thailand, Malaysia, Borneo and Sumatra. It is a small, slender, livebearing fish, with the elongated lower jaw characteristic of its family. The colour of this species varies, depending on where the specimen is found. It is the type species of the genus Dermogenys. Wrestling halfbeaks are surface-feeding fish and feed on a variety of small invertebrates including crustaceans and insect larvae, but especially mosquito larvae and flying insects that have fallen onto the surface of the water. As with all halfbeaks, the upper jaw lifts upwards when the fish is opening its mouth. Wrestling halfbeaks are livebearing fish, the females giving birth to around twenty offspring after a gestation period of about a month. Wrestling halfbeaks are sexually dimorphic. The females are larger than the males and grow up to long; males only reach about and typically has red or yellow patches on the dorsal fin and the beak. The males of wrestling halfbeaks will fight among themselves by locking jaws, hence their name, for up to thirty minutes. This species was described as Dermogenys pusillus by Heinrich Kuhl and Johan Coenraad van Hasselt in 1823 with the type locality given as Bogor, Java, the name was subsequently amended to the feminine form. Fish fighting In the wild, and in large aquaria, the weaker male will quickly disengage and swim away, and fights therefore rarely result in serious injury to either party. However, in their native range, local people sometimes use wrestling halfbeaks as fighting fish for betting purposes (like fighting cocks or siamese fighting fish). Wrestling halfbeaks in aquaria Wrestling halfbeaks, as well as other species in the genus Dermogenys, are quite widely traded as aquarium fish, sometimes under the "silver halfbeak" or "golden halfbeak" names, depending on the colouration of the fish. As with all freshwater halfbeaks, these fish are sensitive to sudden changes in pH and hardness, but they are otherwise adaptable, and can be maintained in anything from soft and acid freshwater through to slightly brackish water. When first introduced into the aquarium, wrestling halfbeaks are nervous fish that tend to be timid. They may swim frantically if suddenly frightened, even crashing into the walls. However, once they are used to their surroundings, they become lively, easy to care for fish.	Biology and health sciences	Acanthomorpha	Animals
1715834	https://en.wikipedia.org/wiki/Water%E2%80%93gas%20shift%20reaction	Water–gas shift reaction	The water–gas shift reaction (WGSR) describes the reaction of carbon monoxide and water vapor to form carbon dioxide and hydrogen: CO + H2O CO2 + H2 The water gas shift reaction was discovered by Italian physicist Felice Fontana in 1780. It was not until much later that the industrial value of this reaction was realized. Before the early 20th century, hydrogen was obtained by reacting steam under high pressure with iron to produce iron oxide and hydrogen. With the development of industrial processes that required hydrogen, such as the Haber–Bosch ammonia synthesis, a less expensive and more efficient method of hydrogen production was needed. As a resolution to this problem, the WGSR was combined with the gasification of coal to produce hydrogen. Applications The WGSR is a highly valuable industrial reaction that is used in the manufacture of ammonia, hydrocarbons, methanol, and hydrogen. Its most important application is in conjunction with the conversion of carbon monoxide from steam reforming of methane or other hydrocarbons in the production of hydrogen. In the Fischer–Tropsch process, the WGSR is one of the most important reactions used to balance the H2/CO ratio. It provides a source of hydrogen at the expense of carbon monoxide, which is important for the production of high purity hydrogen for use in ammonia synthesis. The water–gas shift reaction may be an undesired side reaction in processes involving water and carbon monoxide, e.g. the rhodium-based Monsanto process. The iridium-based Cativa process uses less water, which suppresses this reaction. Fuel cells The WGSR can aid in the efficiency of fuel cells by increasing hydrogen production. The WGSR is considered a critical component in the reduction of carbon monoxide concentrations in cells that are susceptible to carbon monoxide poisoning such as the proton-exchange membrane (PEM) fuel cell. The benefits of this application are two-fold: not only would the water gas shift reaction effectively reduce the concentration of carbon monoxide, but it would also increase the efficiency of the fuel cells by increasing hydrogen production. Unfortunately, current commercial catalysts that are used in industrial water gas shift processes are not compatible with fuel cell applications. With the high demand for clean fuel and the critical role of the water gas shift reaction in hydrogen fuel cells, the development of water gas shift catalysts for the application in fuel cell technology is an area of current research interest. Catalysts for fuel cell application would need to operate at low temperatures. Since the WGSR is slow at lower temperatures where equilibrium favors hydrogen production, WGS reactors require large amounts of catalysts, which increases their cost and size beyond practical application. The commercial LTS catalyst used in large scale industrial plants is also pyrophoric in its inactive state and therefore presents safety concerns for consumer applications. Developing a catalyst that can overcome these limitations is relevant to implementation of a hydrogen economy. Sorption enhanced water gas shift The WGS reaction is used in combination with the solid adsorption of CO2 in the sorption enhanced water gas shift (SEWGS) in order to produce a high pressure hydrogen stream from syngas. Reaction conditions The equilibrium of this reaction shows a significant temperature dependence and the equilibrium constant decreases with an increase in temperature, that is, higher hydrogen formation is observed at lower temperatures. Temperature dependence With increasing temperature, the reaction rate increases, but hydrogen production becomes less favorable thermodynamically since the water gas shift reaction is moderately exothermic; this shift in chemical equilibrium can be explained according to Le Chatelier's principle. Over the temperature range of 600–2000 K, the equilibrium constant for the WGSR has the following relationship: Practical concerns In order to take advantage of both the thermodynamics and kinetics of the reaction, the industrial scale water gas shift reaction is conducted in multiple adiabatic stages consisting of a high temperature shift (HTS) followed by a low temperature shift (LTS) with intersystem cooling. The initial HTS takes advantage of the high reaction rates, but results in incomplete conversion of carbon monoxide. A subsequent low temperature shift reactor lowers the carbon monoxide content to <1%. Commercial HTS catalysts are based on iron oxide–chromium oxide and the LTS catalyst is a copper-based. The copper catalyst is susceptible to poisoning by sulfur. Sulfur compounds are removed prior to the LTS reactor by a guard bed. An important limitation for the HTS is the H2O/CO ratio where low ratios may lead to side reactions such as the formation of metallic iron, methanation, carbon deposition, and the Fischer–Tropsch reaction. High temperature shift catalysis The typical composition of commercial HTS catalyst has been reported as 74.2% Fe2O3, 10.0% Cr2O3, 0.2% MgO (remaining percentage attributed to volatile components). The chromium acts to stabilize the iron oxide and prevents sintering. The operation of HTS catalysts occurs within the temperature range of 310 °C to 450 °C. The temperature increases along the length of the reactor due to the exothermic nature of the reaction. As such, the inlet temperature is maintained at 350 °C to prevent the exit temperature from exceeding 550 °C. Industrial reactors operate at a range from atmospheric pressure to 8375 kPa (82.7 atm). The search for high performance HT WGS catalysts remains an intensive topic of research in fields of chemistry and materials science. Activation energy is a key criteria for the assessment of catalytic performance in WGS reactions. To date, some of the lowest activation energy values have been found for catalysts consisting of copper nanoparticles on ceria support materials, with values as low as Ea = 34 kJ/mol reported relative to hydrogen generation. Low temperature shift catalysis Catalysts for the lower temperature WGS reaction are commonly based on copper or copper oxide loaded ceramic phases, While the most common supports include alumina or alumina with zinc oxide, other supports may include rare earth oxides, spinels or perovskites. A typical composition of a commercial LTS catalyst has been reported as 32-33% CuO, 34-53% ZnO, 15-33% Al2O3. The active catalytic species is CuO. The function of ZnO is to provide structural support as well as prevent the poisoning of copper by sulfur. The Al2O3 prevents dispersion and pellet shrinkage. The LTS shift reactor operates at a range of 200–250 °C. The upper temperature limit is due to the susceptibility of copper to thermal sintering. These lower temperatures also reduce the occurrence of side reactions that are observed in the case of the HTS. Noble metals such as platinum, supported on ceria, have also been used for LTS. Mechanism The WGSR has been extensively studied for over a hundred years. The kinetically relevant mechanism depends on the catalyst composition and the temperature. Two mechanisms have been proposed: an associative Langmuir–Hinshelwood mechanism and a redox mechanism. The redox mechanism is generally regarded as kinetically relevant during the high-temperature WGSR (> 350 °C) over the industrial iron-chromia catalyst. Historically, there has been much more controversy surrounding the mechanism at low temperatures. Recent experimental studies confirm that the associative carboxyl mechanism is the predominant low temperature pathway on metal-oxide-supported transition metal catalysts. Associative mechanism In 1920 Armstrong and Hilditch first proposed the associative mechanism. In this mechanism CO and H2O are adsorbed onto the surface of the catalyst, followed by formation of an intermediate and the desorption of H2 and CO2. In general, H2O dissociates onto the catalyst to yield adsorbed OH and H. The dissociated water reacts with CO to form a carboxyl or formate intermediate. The intermediate subsequently dehydrogenates to yield CO2 and adsorbed H. Two adsorbed H atoms recombine to form H2. There has been significant controversy surrounding the kinetically relevant intermediate during the associative mechanism. Experimental studies indicate that both intermediates contribute to the reaction rate over metal oxide supported transition metal catalysts. However, the carboxyl pathway accounts for about 90% of the total rate owing to the thermodynamic stability of adsorbed formate on the oxide support. The active site for carboxyl formation consists of a metal atom adjacent to an adsorbed hydroxyl. This ensemble is readily formed at the metal-oxide interface and explains the much higher activity of oxide-supported transition metals relative to extended metal surfaces. The turn-over-frequency for the WGSR is proportional to the equilibrium constant of hydroxyl formation, which rationalizes why reducible oxide supports (e.g. CeO2) are more active than irreducible supports (e.g. SiO2) and extended metal surfaces (e.g. Pt). In contrast to the active site for carboxyl formation, formate formation occurs on extended metal surfaces. The formate intermediate can be eliminated during the WGSR by using oxide-supported atomically dispersed transition metal catalysts, further confirming the kinetic dominance of the carboxyl pathway. Redox mechanism The redox mechanism involves a change in the oxidation state of the catalytic material. In this mechanism, CO is oxidized by an O-atom intrinsically belonging to the catalytic material to form CO2. A water molecule undergoes dissociative adsorption at the newly formed O-vacancy to yield two hydroxyls. The hydroxyls disproportionate to yield H2 and return the catalytic surface back to its pre-reaction state. Homogeneous models The mechanism entails nucleophilic attack of water or hydroxide on a M-CO center, generating a metallacarboxylic acid. Thermodynamics The WGSR is exergonic, with the following thermodynamic parameters at room temperature (298 K) {\|class=wikitable !Free energy \|ΔG⊖ = –28.6 kJ/mol \|- !Enthalpy \|ΔH⊖ = –41.2 kJ/mol \|- !Entropy \|ΔS⊖ = –41.84 J/K.mol \|} In aqueous solution, the reaction is less exergonic. Reverse water–gas shift In the conversion of carbon dioxide to useful materials, the water–gas shift reaction is used to produce carbon monoxide from hydrogen and carbon dioxide. This is sometimes called the reverse water–gas shift reaction. Water gas is defined as a fuel gas consisting mainly of carbon monoxide (CO) and hydrogen (H2). The term 'shift' in water–gas shift means changing the water gas composition (CO:H2) ratio. The ratio can be increased by adding CO2 or reduced by adding steam to the reactor.	Physical sciences	Other reactions	Chemistry
15094186	https://en.wikipedia.org/wiki/Graph%20automorphism	Graph automorphism	In the mathematical field of graph theory, an automorphism of a graph is a form of symmetry in which the graph is mapped onto itself while preserving the edge–vertex connectivity. Formally, an automorphism of a graph is a permutation of the vertex set , such that the pair of vertices form an edge if and only if the pair also form an edge. That is, it is a graph isomorphism from to itself. Automorphisms may be defined in this way both for directed graphs and for undirected graphs. The composition of two automorphisms is another automorphism, and the set of automorphisms of a given graph, under the composition operation, forms a group, the automorphism group of the graph. In the opposite direction, by Frucht's theorem, all groups can be represented as the automorphism group of a connected graph – indeed, of a cubic graph. Computational complexity Constructing the automorphism group of a graph, in the form of a list of generators, is polynomial-time equivalent to the graph isomorphism problem, and therefore solvable in quasi-polynomial time, that is with running time for some fixed . Consequently, like the graph isomorphism problem, the problem of finding a graph's automorphism group is known to belong to the complexity class NP, but not known to be in P nor to be NP-complete, and therefore may be NP-intermediate. The easier problem of testing whether a graph has any symmetries (nontrivial automorphisms), known as the graph automorphism problem, also has no known polynomial time solution. There is a polynomial time algorithm for solving the graph automorphism problem for graphs where vertex degrees are bounded by a constant. The graph automorphism problem is polynomial-time many-one reducible to the graph isomorphism problem, but the converse reduction is unknown. By contrast, hardness is known when the automorphisms are constrained in a certain fashion; for instance, determining the existence of a fixed-point-free automorphism (an automorphism that fixes no vertex) is NP-complete, and the problem of counting such automorphisms is ♯P-complete. Algorithms, software and applications While no worst-case polynomial-time algorithms are known for the general Graph Automorphism problem, finding the automorphism group (and printing out an irredundant set of generators) for many large graphs arising in applications is rather easy. Several open-source software tools are available for this task, including NAUTY, BLISS and SAUCY. SAUCY and BLISS are particularly efficient for sparse graphs, e.g., SAUCY processes some graphs with millions of vertices in mere seconds. However, BLISS and NAUTY can also produce Canonical Labeling, whereas SAUCY is currently optimized for solving Graph Automorphism. An important observation is that for a graph on vertices, the automorphism group can be specified by no more than generators, and the above software packages are guaranteed to satisfy this bound as a side-effect of their algorithms (minimal sets of generators are harder to find and are not particularly useful in practice). It also appears that the total support (i.e., the number of vertices moved) of all generators is limited by a linear function of , which is important in runtime analysis of these algorithms. However, this has not been established for a fact, as of March 2012. Practical applications of Graph Automorphism include graph drawing and other visualization tasks, solving structured instances of Boolean Satisfiability arising in the context of Formal verification and Logistics. Molecular symmetry can predict or explain chemical properties. Symmetry display Several graph drawing researchers have investigated algorithms for drawing graphs in such a way that the automorphisms of the graph become visible as symmetries of the drawing. This may be done either by using a method that is not designed around symmetries, but that automatically generates symmetric drawings when possible, or by explicitly identifying symmetries and using them to guide vertex placement in the drawing. It is not always possible to display all symmetries of the graph simultaneously, so it may be necessary to choose which symmetries to display and which to leave unvisualized. Graph families defined by their automorphisms Several families of graphs are defined by having certain types of automorphisms: An asymmetric graph is an undirected graph with only the trivial automorphism. A vertex-transitive graph is an undirected graph in which every vertex may be mapped by an automorphism into any other vertex. An edge-transitive graph is an undirected graph in which every edge may be mapped by an automorphism into any other edge. A symmetric graph is a graph such that every pair of adjacent vertices may be mapped by an automorphism into any other pair of adjacent vertices. A distance-transitive graph is a graph such that every pair of vertices may be mapped by an automorphism into any other pair of vertices that are the same distance apart. A semi-symmetric graph is a graph that is edge-transitive but not vertex-transitive. A half-transitive graph is a graph that is vertex-transitive and edge-transitive but not symmetric. A skew-symmetric graph is a directed graph together with a permutation σ on the vertices that maps edges to edges but reverses the direction of each edge. Additionally, σ is required to be an involution. Inclusion relationships between these families are indicated by the following table:	Mathematics	Graph theory	null
2389770	https://en.wikipedia.org/wiki/Panthera%20hybrid	Panthera hybrid	A Panthera hybrid is a crossbreed between individuals of any of the five species of the genus Panthera: the tiger, lion, jaguar, leopard, and snow leopard. Most hybrids would not be perpetuated in the wild as the territories of the parental species do not overlap and the males are usually infertile. Mitochondrial genome research revealed that wild hybrids were also present in ancient times. The mitochondrial genomes of the snow leopard and the lion were more similar to each other than to other Panthera species, indicating that at some point in their history, the female hybrid progeny of male ancestors of modern snow leopards and female ancestors of modern lions interbred with male ancestors of modern snow leopards. History In theory, lions and tigers can be matched in the wild and give offspring. In reality, there may be no natural born tigon or liger in the world, as lions and tigers are separated both geographically and by behavioral differences. In England, African lions and Asian tigresses have been successfully mated, and three lion-tiger hybrid cubs were born in Windsor in 1824, which is probably the earliest record of captive-bred ligers. The three cubs were then presented to George IV. Naming of hybrids Panthera hybrids are typically given a portmanteau name, varying by which species is the sire (male parent) and which is the dam (female parent). For example, a hybrid between a lion and a tigress is a liger, because the lion is the male parent and the tigress is the female parent. Jaguar and leopard hybrids A jagupard, jagulep or jagleop is the hybrid of a male jaguar and a female leopard. A single rosetted female jagupard was produced at a zoo in Chicago, United States. Jaguar-leopard hybrids bred at Hellbrun Zoo, Salzburg were described as jagupards, which conforms to the usual portmanteau naming convention. A leguar or lepjag is the hybrid of a male leopard and a female jaguar. The terms jagulep and lepjag are often used interchangeably, regardless of which animal was the sire. Numerous lepjags have been bred as animal actors, as they are more tractable than jaguars. The 19th century zoologist A.D. Bartlett stated: "I have, more than once, met with instances of the male jaguar (P. onca) breeding with a female leopard (P. pardus). These hybrids were also reared recently in Wombell's well-known travelling collection. I have seen some animals of this kind bred, between a male black jaguar and a female Indian leopard:-the young partook strongly of the male, being almost black." In Barnabos Menagerie (in Spain), a jaguar gave birth to two cubs from a union with a black leopard; one resembled the dam, but was somewhat darker, while the other was black with the rosettes of the dam showing. Since melanism in the panther (leopard) is recessive, the jaguar would have had to have been black, or be a jaguar-black leopard hybrid itself, carrying the recessive gene. Scherren continued, "The same cross, but with the sexes reversed, was noted, by Professor Sacc (F) of Barcelona Zoo (Zoolog. Gart., 1863, 88). "The cub, a female, was grey. She is said to have produced two cubs to her sire; one like a jaguar, the other like the dam. Herr Rorig expressed his regret that the account of the last two cases mentioned lacked fullness and precision." Female jaguleps or lepjags are fertile, and when one is mated to a male lion, the offspring are referred to as lijaguleps. One such complex hybrid was exhibited in the early 1900s as a "Congolese spotted lion", hinting at some exotic African beast, rather than a man-made hybrid. Jaguar and lion hybrids A jaglion or jaguon is the offspring between a male jaguar and a female lion (lioness). A mounted specimen is on display at the Walter Rothschild Zoological Museum, Hertfordshire, England. It has the lion's background color, brown, jaguar-like rosettes and the powerful build of the jaguar. On April 9, 2006, two jaglions were born at Bear Creek Wildlife Sanctuary, Barrie (north of Toronto), Ontario, Canada. Jahzara (female) and Tsunami (male) were the result of an unintended mating between a black jaguar called Diablo and a lioness called Lola, which had been hand-raised together and were inseparable. They were kept apart when Lola came into oestrus. Tsunami is spotted, but Jahzara is a melanistic jaglion due to inheriting the jaguar's dominant melanism gene. It was not previously known how the jaguar's dominant melanism gene would interact with lion coloration genes. A liguar is an offspring of a male lion and a female jaguar. When the fertile offspring of a male lion and female jaguar mates with a leopard, the resulting offspring is referred to as a leoliguar. Jaguar and tiger hybrids A tiguar is an offspring of a male tiger and a female jaguar. Reportedly, at the Altiplano Zoo in the city of San Pablo Apetatitlán (near Tlaxcala City, Mexico), the crossbreeding of a male Siberian tiger and a female jaguar from the southern Lacandon Jungle produced a male tiguar named Mickey. Mickey was on exhibition at a 400 m2 habitat and as of June 2009, was two years old and weighed . Attempts to verify this report have been bolstered by recent images purported to show the adult Mickey (see External links section). There has been no report of the birth of a hybrid from a male jaguar and female tiger, which would be termed a "jagger". There is a claimed sighting of a lion × black jaguar cross (male) and a tiger × black jaguar cross (female) loose in Maui, Hawaii. There are no authenticated tiger/jaguar hybrids and the description matches that of a liger. The alleged tiger × black jaguar was large, relatively long-necked (probably due to lack of a ruff or mane) with both stripes and "jaguar-like" rosettes on its sides. The assertion of hybrid identity was due to the combination of black, dark brown, light brown, dark orange, dark yellow and beige markings and the tiger-like stripes radiating from its face. It is more likely to have been a released liger, since these are very large and have a mix of rosettes (lion juvenile markings) and stripes and can have a brindled mix of colors exactly as described (their markings are extremely variable). Leopard and lion hybrids A leopon is the result of breeding a leopard and a lioness. They occur only in captivity. The first documented leopon was bred at Kolhapur, India, in 1910. Its skin was sent to Reginald Innes Pocock by Walter Samuel Millard, the Secretary of the Bombay Natural History Society. It was a cross between a large leopard and a lioness. Two cubs were born, one of which died aged 2.5 months, and the other was still living when Pocock described it in 1912. Pocock wrote that it was spotted like a leopard, but that the spots on its sides were smaller and closer set than those of an Indian leopard and were brown and indistinct, like the fading spots of a juvenile lion. The spots on the head, spine, belly and legs were black and distinct. The tail was spotted on the topside and striped underneath and had a blackish tip with longer hairs. The underside was dirty white, the ears were fawn and had a broad black bar, but did not have the white spot found in leopards. Another lion-leopardess hybrid was born in Florence, Italy called lionard or lipard (/'laɪpəd/ or /'laɪpərd/). Leopard and tiger hybrids The name dogla is a native Indian name used for a supposedly natural hybrid offspring of a male leopard and a female tiger (tigress). Indian folklore claims that large male leopards sometimes mate with tigresses, and anecdotal evidence exists in India of offspring resulting from leopard to tigress matings. A supposed dogla was reported in the early 1900s. Tiger-leopardess hybrids have supposedly appeared many times. Frederick Codrington Hicks recorded that the weight of these creatures varied from 50 pounds to the weight of a tigress. In addition, in September 1965, a "leoger" skin was supposedly put on sale. There are some more documentations of this hybrid, but most of them are just of strange-looking skins that could also be attributed to genetic mutations. Most of these reports are probably hoaxes or misinterpretations, which makes it hard for scientists to learn about tiger-leopardess hybrids, but at least a part of the claims are true or in part true, such as the ones made by Frederick Codrington Hicks. K Sankhala's book Tiger refers to large, troublesome leopards as adhabaghera, which he translated as "bastard", and suggests a leopard/tiger hybrid (the reverse hybrid is unlikely to arise in the wild state, as a wild male tiger would probably kill rather than mate with a female leopard). Sankhala noted there was a belief amongst local people that leopards and tigers naturally hybridise. From "The Tiger, Symbol Of Freedom", edited by Nicholas Courtney: "Rare reports have been made of tigresses mating with leopards in the wild. There has even been an account of the sighting of rosettes; the stripes of the tiger being most prominent in the body. The animal was a male measuring a little over eight feet [2.44 m]." This is the same description as given by Hicks. The 1951 book Mammalian Hybrids reported tiger/leopard matings were infertile, producing spontaneously aborted "walnut-sized fetuses". A tigard is the hybrid offspring of a tiger and a leopardess. The only known attempts to mate the two have produced stillborns. In 1900, Carl Hagenbeck crossed a female leopard with a Bengal tiger. The stillborn offspring had a mixture of spots, rosettes and stripes. Henry Scherren wrote, "A male tiger from Penang served two female Indian leopards, and twice with success. Details are not given and the story concludes somewhat lamely. 'The leopardess dropped her cubs prematurely, the embryos were in the first stage of development and were scarcely as big as young mice.' Of the second leopardess there is no mention." Lion and tiger hybrids The hybrids resulting from crossbreeding between lions and tigers are known as tigon (/ˈtaɪɡən/) and liger (/ˈlaɪɡər/). The second generation hybrids of liger or tigon are known as liliger, tiliger, litigon and titigon. The tigon (Panthera tigris X leo), also known as tiglon (/ˈtaɪɡlən/) is an offspring of a male tiger (Panthera tigris) and a female lion (Panthera leo). A liger is distinct from tigon (Panthera leo X tigris), as a hybrid of female tiger and male lion. Professor Valentine Bail conducted a long observation and recording of some lion-tiger hybrids, those lion-tiger are owned by Mr. Atkins and his zoo: The early record lion-tiger hybrid was mainly tigons, in At Home In The Zoo (1961), Gerald Iles wrote "For the record I must say that I have never seen a liger, a hybrid obtained by crossing a lion with a tigress. They seem to be even rarer than tigons." Liger A liger is the offspring between a male lion and a female tiger, which is larger than its parents because the lion has a growth maximizing gene and the tigress, unlike the lioness, has no growth inhibiting gene. Tigon A tigon is the offspring of a female lion and a male tiger. The tigon is not as common as the converse hybrid, the liger. Contrary to some beliefs, the tigon ends up smaller than either parent, because male tigers and lionesses have a growth inhibitor. In the late 19th and early 20th centuries, tigons were more common than ligers. Liliger A liliger is the offspring of a lion and a ligress. The first known liliger is a cub named Kiara. Litigon Rudrani, a tigoness from the Alipore Zoo, mated with Debabrata, a male lion, and gave birth to three litigons. Only one litigon cub, named Cubanacan, survived. Tiliger A tiliger is the offspring of a male tiger and a ligress. Titigon A titigon is the offspring of a male tiger and a tigoness. Growth and size Typically, ligers are more likely to be larger and heavier than all other existing felids. Some biologists believe that their gigantism results from the lack of certain genes that limit the growth of lions. The male lion's genes tend to maximize the growth of its progeny, as the larger size represents greater competitiveness. In order to control the size of the offspring within a certain range, the growth-inhibiting gene of the lioness will offset the growth-maximizing gene of the male lion. The genes of a female tiger, however, are not adapted to limiting growth, which allows ligers to grow far larger and heavier than either parent. In general, most ligers grow more than in length and weigh more than . According to the Guinness World Records (through 2013), the world's largest felid was the adult male liger, Hercules, from Myrtle Beach Safari, a wildlife reserve in South Carolina, US. He was measured at (standing at the shoulder) and weighed at . Hercules eats approximately of meat and drinks several liters of water per day. Tigons are a cross between a male tiger and a female lion. The presence of growth-minimizing genes from the lioness causes them to be smaller than either of their parent species; they weigh less than . Tigons also have growth dysplasia (however, inversely). A tigon is approximately twice as light as liger. Appearance Ligers and tigons look similar to their parents, only bigger or smaller. Their teeth are about two inches long. They have the genetic components of tigers and lions; therefore, they may be very similar to their parent species and can be difficult to identify. Their coloring ranges from gold to brown to white, and they may have spots or stripes. Adult male ligers usually have smaller manes than male lions. Longevity A liger called Samson died at the age of thirteen in 2006. Shasta, a female liger, was born in the Hogle Zoo in Salt Lake City in 1948, and died in 1972. She lived for 24 years. Many claim that ligers are short-lived, but according to the survey, such a conclusion is still uncertain. A male tigon owned by Atkins born on July 19, 1833, lived for 10 years. Fertility Guggisberg said ligers and tigons were thought to be invariably sterile. The first hybrid of a hybrid, a cub mothered by a liger, was discovered at the Munich-Hellabrunn Zoo in 1943. The birth of a second generation of hybrids proved that the biologists were wrong about tigons' and ligers' fertility; it now seems that only male lion-tiger hybrids are sterile. Zoo animals By 2017, roughly more than 100 ligers were thought to exist; however, only a few tigons still exist, as they are more difficult to breed. Moreover, ligers are more likely to attract tourists, so zoos prefer to breed ligers as opposed to tigons. Some zoos claim they breed ligers or tigons for conservation, but opponents believe that it is meaningless to preserve a species that does not exist in the wild.	Biology and health sciences	Hybrids	Animals
2390955	https://en.wikipedia.org/wiki/Youngina	Youngina	Youngina (named after John Young (1823–1900)) is an extinct genus of diapsid reptile from the Late Permian Beaufort Group (Tropidostoma-Dicynodon zones) of the Karoo Red Beds of South Africa. This, and a few related forms, make up the family Younginidae, within the order Eosuchia (proposed by Broom in 1914). Eosuchia, having become a wastebasket taxon for many probably distantly-related primitive diapsid reptiles ranging from the Late Carboniferous to the Eocene, Romer proposed that it be replaced by Younginiformes (that included Younginidae and the Tangasauridae, ranging from the Permian to the Triassic). Taxonomy Youngina is known from several specimens. Many of these were attributed to as separate genera and species (such as Youngoides and Youngopsis), but it was later realized that they were not distinct from Y. capensis. The holotype specimen of Youngina, discovered by Broom himself, was described briefly in 1914. The "Youngoides romeri" specimen was first attributed to Youngina, but later given its eponymous and separate designation in a later paper. Acanthotoposaurus is also a junior synonym of Youngina. Description Youngina was a relatively small reptile, with a skull length of and a total body length of . The braincase anatomy was redescribed in 2010. Youngina shows a mosaic of features found in more primitive diapsids and more derived taxa such as archosauromorphs and lepidosauromorphs suggesting a non-orthogenetic evolution of these characters. Though the palatobasal articulation is open, it was probably immobile, similar to the skull of the tuatara, contrary to some earlier claims made about the metakinetic mobility of basicranial joints in Youngina and other early diapsid reptiles. Phylogeny Youngina was once thought to be closely related to Acerosodontosaurus, and more distantly to tangasaurids (Kenyasaurus, Hovasaurus, Thadeosaurus, and Tangasaurus), but the monophyly of has not been demonstrated in published analyses of diapsid reptiles, and it is likely this group is paraphyletic. Acerosodontosaurus is probably closer to other former , rather than being closely related to Youngina. Below is a cladogram from the analysis of Reisz et al. (2011) showing the phylogenetic position of Youngina among early diapsids:	Biology and health sciences	Other prehistoric reptiles	Animals
2391004	https://en.wikipedia.org/wiki/Ribbonfish	Ribbonfish	The ribbonfish are any lampriform fishes in the family Trachipteridae. There are about 10 recognized species in the family. These pelagic fish are named for their slim, ribbon-like appearance. They are rarely seen alive, as they typically live in deep waters, though are not bottom feeders. The perciform fish known as the red bandfish (Cepola macrophthalma) is sometimes referred to as ribbonfish, but it is unrelated to any ribbonfish in the Trachipteridae. They are readily recognized by their anatomy — a long, compressed, tape-like body, short head, narrow mouth and feeble teeth. A high dorsal fin occupies the whole length of the back; an anal fin is absent, and the caudal fin, if present, consists of two fascicles of rays of which the upper is prolonged and directed upwards. The pectoral fins are small, the pelvic fins composed of several rays, or of one long ray only. They have heavy spines along their lateral lines, and numerous lumps in the skin. Ribbonfish possess all the characteristics of fish living at very great depths. Their fins especially, and the membrane connecting them, are of a very delicate and brittle structure. In young ribbonfish, some of the fin-rays are prolonged to an extraordinary degree, and sometimes provided with appendages. Specimens have been taken in the Atlantic, the Mediterranean, the Bay of Bengal, at Mauritius, and in the Pacific. The species from the Atlantic has occurred chiefly on the northern coasts, Iceland, Scandinavia, Orkney, and Scotland. The north Atlantic species is known in English as deal fish, in Icelandic as vogmær and in Swedish as vågmär. Its length is usually 5 to 8 ft (1.5–3.5 m), but it can sometimes be found at over 20 ft. Specimens seem usually to be driven to the shore by gales in winter, and are sometimes left by the tide. S. Nilsson, however, in Scandinavia observed a living specimen in two or three fathoms (4–5 m) of water moving something like a flatfish with one side turned obliquely upwards. A specimen of Trachipterus ishikawae was discovered on a beach in Kenting, Taiwan, in November 2007, alive but with a 10-cm cut wound to its side, and was returned to deeper water. The species Trachipterus ishikawae is commonly called "earthquake fish" in Taiwan because the fish are popularly believed to appear following major earthquake events due to alleged sensitivity to disturbances in the ocean floor. Records of such appearances were made following a 100-year earthquake in Hengchun in late 2006 and in Taitung in 2007, as well as the numerous March 2011 sightings along the coast of Japan, but other recorded sightings do not correspond with seismic disturbances.	Biology and health sciences	Acanthomorpha	Animals
2391490	https://en.wikipedia.org/wiki/Pig	Pig	The pig (Sus domesticus), also called swine (: swine) or hog, is an omnivorous, domesticated, even-toed, hoofed mammal. It is named the domestic pig when distinguishing it from other members of the genus Sus. It is considered a subspecies of Sus scrofa (the wild boar or Eurasian boar) by some authorities, but as a distinct species by others. Pigs were domesticated in the Neolithic, both in East Asia and in the Near East. When domesticated pigs arrived in Europe, they extensively interbred with wild boar but retained their domesticated features. Pigs are farmed primarily for meat, called pork. The animal's skin or hide is used for leather. China is the world's largest pork producer, followed by the European Union and then the United States. Around 1.5 billion pigs are raised each year, producing some 120 million tonnes of meat, often cured as bacon. Some are kept as pets. Pigs have featured in human culture since Neolithic times, appearing in art and literature for children and adults, and celebrated in cities such as Bologna for their meat products. Description The pig has a large head, with a long snout strengthened by a special prenasal bone and a disk of cartilage at the tip. The snout is used to dig into the soil to find food and is an acute sense organ. The dental formula of adult pigs is , giving a total of 44 teeth. The rear teeth are adapted for crushing. In males, the canine teeth can form tusks, which grow continuously and are sharpened by grinding against each other. There are four hoofed toes on each foot; the two larger central toes bear most of the weight, while the outer two are also used in soft ground. Most pigs have rather sparsely bristled hair on their skin, though there are some woolly-coated breeds such as the Mangalitsa. Adult pigs generally weigh between , though some breeds can exceed this range. Exceptionally, a pig called Big Bill weighed and had a shoulder height of . Pigs possess both apocrine and eccrine sweat glands, although the latter are limited to the snout. Pigs, like other "hairless" mammals such as elephants, do not use thermal sweat glands in cooling. Pigs are less able than many other mammals to dissipate heat from wet mucous membranes in the mouth by panting. Their thermoneutral zone is . At higher temperatures, pigs lose heat by wallowing in mud or water via evaporative cooling, although it has been suggested that wallowing may serve other functions, such as protection from sunburn, ecto-parasite control, and scent-marking. Pigs are among four mammalian species with mutations in the nicotinic acetylcholine receptor that protect against snake venom. Mongooses, honey badgers, hedgehogs, and pigs all have different modifications to the receptor pocket which prevents α-neurotoxin from binding. Pigs have small lungs for their body size, and are thus more susceptible than other domesticated animals to fatal bronchitis and pneumonia. The genome of the pig has been sequenced; it contains about 22,342 protein-coding genes. Evolution Phylogeny Domestic pigs are related to other pig species as shown in the cladogram, based on phylogenetic analysis using mitochondrial DNA. Taxonomy The pig is most often considered to be a subspecies of the wild boar, which was given the name Sus scrofa by Carl Linnaeus in 1758; following from this, the formal name of the pig is Sus scrofa domesticus. However, in 1777, Johann Christian Polycarp Erxleben classified the pig as a separate species from the wild boar. He gave it the name Sus domesticus, still used by some taxonomists. The American Society of Mammalogists considers it a separate species. Domestication in the Neolithic Archaeological evidence shows that pigs were domesticated from wild boar in the Near East in or around the Tigris Basin, being managed in a semi-wild state much as they are managed by some modern New Guineans. There were pigs in Cyprus more than 11,400 years ago, introduced from the mainland, implying domestication in the adjacent mainland by then. Pigs were separately domesticated in China, starting some 8,000 years ago. In the Near East, pig husbandry spread for the next few millennia. It reduced gradually during the Bronze Age, as rural populations instead focused on commodity-producing livestock, but it was sustained in cities. Domestication did not involve reproductive isolation with population bottlenecks. Western Asian pigs were introduced into Europe, where they crossed with wild boar. There appears to have been interbreeding with a now extinct ghost population of wild pigs during the Pleistocene. The genomes of domestic pigs show strong selection for genes affecting behavior and morphology. Human selection for domestic traits likely counteracted the homogenizing effect of gene flow from wild boars and created domestication islands in the genome. Pigs arrived in Europe from the Near East at least 8,500 years ago. Over the next 3,000 years they interbred with European wild boar until their genome showed less than 5% Near Eastern ancestry, yet retained their domesticated features. DNA evidence from subfossil remains of teeth and jawbones of Neolithic pigs shows that the first domestic pigs in Europe were brought from the Near East. This stimulated the domestication of local European wild boar, resulting in a third domestication event with the Near Eastern genes dying out in European pig stock. More recently there have been complex exchanges, with European domesticated lines being exported, in turn, to the ancient Near East. Historical records indicate that Asian pigs were again introduced into Europe during the 18th and early 19th centuries. History Columbian Exchange Among the animals that the Spanish introduced to the Chiloé Archipelago in the 16th century Columbian Exchange, pigs were the most successful in adapting to local conditions. The pigs benefited from abundant shellfish and algae exposed by the large tides of the archipelago. Pigs were brought to southeastern North America from Europe by de Soto and other early Spanish explorers. Escaped pigs became feral. Feral pigs Pigs have escaped from farms and gone feral in many parts of the world. Feral pigs in the southeastern United States have migrated north to the Midwest, where many state agencies have programs to remove them. Feral pigs in New Zealand and northern Queensland have caused substantial environmental damage. Feral hybrids of the European wild boar with the domestic pig are disruptive to both environment and agriculture, as they destroy crops, spread animal diseases including Foot-and-mouth disease, and consume wildlife such as juvenile seabirds and young tortoises. Feral pig damage is especially an issue in southeastern South America. Reproduction Physiology Female pigs reach sexual maturity at 3–12 months of age and come into estrus every 18–24 days if they are not successfully bred. The variation in ovulation rate can be attributed to intrinsic factors such as age and genotype, as well as extrinsic factors like nutrition, environment, and the supplementation of exogenous hormones. The gestation period averages 112–120 days. Estrus lasts two to three days, and the female's displayed receptiveness to mate is known as standing heat. Standing heat is a reflexive response that is stimulated when the female is in contact with the saliva of a sexually mature boar. Androstenol is one of the pheromones produced in the submaxillary salivary glands of boars that trigger the female's response. The female cervix contains a series of five interdigitating pads, or folds, that hold the boar's corkscrew-shaped penis during copulation. Females have bicornuate uteruses and two conceptuses must be present in both uterine horns to enable pregnancy to proceed. The mother's body recognises that it is pregnant on days 11 to 12 of pregnancy, and is marked by the corpus luteum's producing the sex hormone progesterone. To sustain the pregnancy, the embryo signals to the corpus luteum with the hormones estradiol and prostaglandin E2. This signaling acts on both the endometrium and luteal tissue to prevent the regression of the corpus luteum by activation of genes that are responsible for corpus luteum maintenance. During mid to late pregnancy, the corpus luteum relies primarily on luteinizing hormone for maintenance until birth. Archeological evidence indicates that medieval European pigs farrowed, or bore a litter of piglets, once per year. By the nineteenth century, European piglets routinely double-farrowed, or bore two litters of piglets per year. It is unclear when this shift occurred. Pigs have a maximum life span of about 27 years. Nest-building A characteristic of pigs which they share with carnivores is nest-building. Sows root in the ground to create depressions the size of their body, and then build nest mounds, using twigs and leaves, softer in the middle, in which to give birth. When the mound reaches the desired height, she places large branches, up to 2 metres in length, on the surface. She enters the mound and roots around to create a depression within the gathered material. She then gives birth in a lying position, unlike other artiodactyls which usually stand while birthing. Nest-building occurs during the last 24 hours before the onset of farrowing, and becomes most intense 12 to 6 hours before farrowing. The sow separates from the group and seeks a suitable nest site with well-drained soil and shelter from rain and wind. This provides the offspring with shelter, comfort, and thermoregulation. The nest provides protection against weather and predators, while keeping the piglets close to the sow and away from the rest of the herd. This ensures they do not get trampled on, and prevents other piglets from stealing milk from the sow. The onset of nest-building is triggered by a rise in prolactin level, caused by a decrease in progesterone and an increase in prostaglandin; the gathering of nest material seems to be regulated more by external stimuli such as temperature. Nursing and suckling Pigs have complex nursing and suckling behaviour. Nursing occurs every 50–60 minutes, and the sow requires stimulation from piglets before milk let-down. Sensory inputs (vocalisation, odours from mammary and birth fluids, and hair patterns of the sow) are particularly important immediately post-birth to facilitate teat location by the piglets. Initially, the piglets compete for position at the udder; then the piglets massage around their respective teats with their snouts, during which time the sow grunts at slow, regular intervals. Each series of grunts varies in frequency, tone and magnitude, indicating the stages of nursing to the piglets. The phase of competition for teats and of nosing the udder lasts for about a minute, ending when milk begins to flow. The piglets then hold the teats in their mouths and suck with slow mouth movements (one per second), and the rate of the sow's grunting increases for approximately 20 seconds. The grunt peak in the third phase of suckling does not coincide with milk ejection, but rather the release of oxytocin from the pituitary into the bloodstream. Phase four coincides with the period of main milk flow (10–20 seconds) when the piglets suddenly withdraw slightly from the udder and start sucking with rapid mouth movements of about three per second. The sow grunts rapidly, lower in tone and often in quick runs of three or four, during this phase. Finally, the flow stops and so does the grunting of the sow. The piglets may dart from teat to teat and recommence suckling with slow movements, or nosing the udder. Piglets massage and suckle the sow's teats after milk flow ceases as a way of letting the sow know their nutritional status. This helps her to regulate the amount of milk released from that teat in future sucklings. The more intense the post-feed massaging of a teat, the more milk that teat later releases. Teat order In pigs, dominance hierarchies are formed at an early age. Piglets are precocious, and attempt to suckle soon after being born. The piglets are born with sharp teeth and fight for the anterior teats, as these produce more milk. Once established, this teat order remains stable; each piglet tends to feed on a particular teat or group of teats. Stimulation of the anterior teats appears to be important in causing milk letdown, so it might be advantageous to the entire litter to have these teats occupied by healthy piglets. Piglets locate teats by sight and then by olfaction. Behaviour Social Pig behaviour is intermediate between that of other artiodactyls and of carnivores. Pigs seek out the company of other pigs and often huddle to maintain physical contact, but they do not naturally form large herds. They live in groups of about 8–10 adult sows, some young individuals, and some single males. Pigs confined in a simplified, crowded, or uncomfortable environment may resort to tail-biting; farmers sometimes dock the tails of pigs to prevent the problem, or may enrich the environment with toys or other objects to reduce the risk. Temperature control Because of their relative lack of sweat glands, pigs often control their body temperature using behavioural thermoregulation. Wallowing, coating the body with mud, is a common behaviour. They do not submerge completely under the mud, but vary the depth and duration of wallowing depending on environmental conditions. Adult pigs start wallowing once the ambient temperature is around . They cover themselves in mud from head to tail. They may use mud as a sunscreen, or to keep parasites away. Most bristled pigs "blow their coat", meaning that they shed most of the longer, coarser stiff hair once a year, usually in spring or early summer, to prepare for the warmer months ahead. Eating, feeding, sleeping Where pigs are allowed to roam freely, they walk roughly 4 km daily, scavenging within a home range of around a hectare. Farmers in Africa often choose such a low-input, free-range production system. If conditions permit, pigs feed continuously for many hours and then sleep for many hours, in contrast to ruminants, which tend to feed for a short time and then sleep for a short time. Pigs are omnivorous and versatile in their feeding behaviour. They primarily eat leaves, stems, roots, fruits, and flowers. Rooting is an instinctual comforting behaviour in pigs characterized by nudging the snout into something. It first happens when piglets are born to obtain their mother's milk, and can become a habitual, obsessive behaviour, most prominent in animals weaned too early. Pigs root and dig into the ground to forage for food. Rooting is also a means of communication. Intelligence Pigs are relatively intelligent animals, roughly on par with dogs. They distinguish each other as individuals, spend time in play, and form structured communities. They have good long-term memory and they experience emotions, changing their behaviour in response to the emotional states of other pigs. In terms of experimental tasks, pigs can perform tasks that require them to identify the locations of objects; they can solve mazes; and they can work with a simple language of symbols. They display self-recognition in a mirror. Pigs have been trained to associate different sorts of music (Bach and a military march) with food and social isolation respectively, and could communicate the resulting positive or negative emotion to untrained pigs. Pigs can be trained to use a joystick with their snout to select a target on screen. Senses Pigs have panoramic vision of approximately 310° and binocular vision of 35° to 50°. It is thought they have no eye accommodation. Other animals that have no accommodation, e.g. sheep, lift their heads to see distant objects. The extent to which pigs have colour vision is still a source of some debate; however, the presence of cone cells in the retina with two distinct wavelength sensitivities (blue and green) suggests that at least some colour vision is present. Pigs have a well-developed sense of smell; this is exploited in Europe where trained pigs find underground truffles. Pigs have 1,113 genes for smell receptors, compared to 1,094 in dogs; this may indicate an acute sense of smell, but against this, insects have only around 50 to 100 such genes but make extensive use of olfaction. Olfactory rather than visual stimuli are used in the identification of other pigs. Hearing is well developed; sounds are localised by moving the head. Pigs use auditory stimuli extensively for communication in all social activities. Alarm or aversive stimuli are transmitted to other pigs not only by auditory cues but also by pheromones. Similarly, recognition between the sow and her piglets is by olfactory and vocal cues. Pests and diseases Pigs are subject to many pests and diseases which can seriously affect productivity and cause death. These include parasites such as Ascaris roundworms, virus diseases such as the tick-borne African Swine Fever, bacterial infections such as Clostridium, arthritis caused by Mycoplasma, and stillbirths caused by Parvovirus. Some parasites of pigs are a public health risk as they can be transmitted to humans in undercooked pork. These are the pork tapeworm Taenia solium; a protozoan, Toxoplasma gondii; and a nematode, Trichinella spiralis. Transmission can be prevented by thorough sanitation on the farm; by meat inspection and careful commercial processing; and by thorough cooking, or alternatively by sufficient freezing and curing. In agriculture Production Pigs have been raised outdoors, and sometimes allowed to forage in woods or pastures. In industrialized nations, pig production has largely switched to large-scale intensive pig farming. This has lowered production costs but has caused concern about possible cruelty. As consumers have become concerned with the humane treatment of livestock, demand for pasture-raised pork in these nations has increased. Most pigs in the US receive ractopamine, a beta-agonist drug, which promotes muscle instead of fat and quicker weight gain, requiring less feed to reach finishing weight, and producing less manure. China has requested that pork exports be ractopamine-free. With a population of around 1 billion individuals, the domesticated pig is one of the most numerous large mammals on the planet. Like all animals, pigs are susceptible to adverse impacts from climate change, such as heat stress from increased annual temperatures and more intense heatwaves. Heat stress has increased rapidly between 1981 and 2017 on pig farms in Europe. Installing a ground-coupled heat exchanger is an effective intervention. Breeds Around 600 breeds of pig have been created by farmers around the world, mainly in Europe and Asia, differing in coloration, shape, and size. According to The Livestock Conservancy, as of 2016, three breeds of pig are critically rare (having a global population of fewer than 2000). They are the Choctaw hog, the Mulefoot, and the Ossabaw Island hog. The smallest known pig breed in the world is the Göttingen minipig, typically weighing about as a healthy, full-grown adult. As pets Vietnamese Pot-bellied pigs, a miniature breed of pig, have been kept as pets in the United States, beginning in the latter half of the 20th century. Pigs are intelligent, social creatures. They are considered hypoallergenic and are known to do quite well with people who have the usual animal allergies. Since these animals are known to have a life expectancy of 15 to 20 years, they require a long-term commitment. Given pigs are bred primarily as livestock and have not been bred as companion animals for very long, selective breeding for a placid or biddable temperament is not well established. Pigs have radically different psychology and behaviours compared to dogs, and exhibit fight-or-flight instincts, an independent nature, and natural assertiveness. Male and female swine that have not been de-sexed may express unwanted aggressive behavior, and are prone to developing serious health issues. As rooting is found to be comforting, pigs kept in the house may root household objects, furniture or surfaces. Pet pigs should be let outside to allow them to fulfill their natural desire of rooting around. Economy Approximately 1.5 billion pigs are slaughtered each year for meat. The pork belly futures contract became an icon of commodities trading. It appears in depictions of the arena in popular entertainment, such as the 1983 film Trading Places. Trade in pork bellies declined, and they were delisted from the Chicago Mercantile Exchange in 2011. In 2023, China produced more pork than any other country, 55 million tonnes, followed by the European Union with 22.8 million tonnes and the United States with 12.5 million tonnes. Global production in 2023 was 120 million tonnes. India, despite its large population, consumed under 0.3 million tonnes of pork in 2023. International trade in pork (meat not consumed in the producing country) reached 13 million tonnes in 2020. Uses Products Pigs are farmed primarily for meat, called pork. Pork is eaten in the form of pork chops, loin or rib roasts, shoulder joints, steaks, and loin (also called fillet). The many meat products made from pork include ham, bacon (mainly from the back and belly), and sausages. Pork is further made into charcuterie products such as terrines, galantines, pâtés and confits. Some sausages such as salami are fermented and air-dried, to be eaten raw. There are many types, the original Italian varieties including Genovese, Milanese, and Cacciatorino, with spicier kinds from the South of Italy including Calabrese, Napoletano, and Peperone. The hide is made into pigskin leather, which is soft and durable; it can be brushed to form suede leather. These are used for products such as gloves, wallets, suede shoes, and leather jackets. In the 16th century, pig skin was the most popular book-binding material in Germany, though calf skin was more common elsewhere. In medicine Pigs, both as live animals and as a source of post-mortem tissues, are valuable animal models because of their biological, physiological, and anatomical similarities to human beings. For instance, human skin is very similar to the pigskin, therefore pigskin has been used in many preclinical studies. Pigs are good non-human candidates for organ donation to humans, and in 2021 became the first animal to successfully donate an organ to a human body. The procedure used a donor pig genetically engineered not to have a specific carbohydrate that the human body considers a threat–Galactose-alpha-1,3-galactose. Pigs are good for human donation as the risk of cross-species disease transmission is reduced by the considerable phylogenetic distance from humans. They are readily available, and the danger of creating new human diseases is low as domesticated pigs have been in close contact with humans for thousands of years. Impact of pig husbandry On public health Pig farms can serve as reservoirs of viral diseases that are dangerous to humans and so contribute to their outbreaks in human populations. The 2009 swine flu pandemic was caused by an influenza A variant which had first emerged in pigs. Pigs were also essential to the first outbreak of the Nipah virus in 1999, with 93% of the infected humans having had contact with pigs. While Japanese encephalitis is primarily spread by mosquitoes, pigs are a known intermediary host. There is also a potential for porcine coronaviruses such as porcine epidemic diarrhea virus or swine acute diarrhea syndrome coronavirus to spill over into human populations. On the environment As with the other forms of meat, producing pork is more energy-intensive than plant-based foods, and it is associated with more greenhouse gas emissions per calorie. However, emissions from pork are many times smaller than those of beef, veal and mutton, though larger than of chicken meat. Intensive pig production is also associated with water pollution concerns, as the swine waste is often stored above ground in so-called lagoons. These lagoons typically have high levels of nitrogen and phosphorus, and can contain toxic heavy metals like zinc and copper, microbial pathogens, or hold elevated concentrations of pharmaceuticals from subtherapeutic antibiotic use in swine. This wastewater from lagoons is liable to reach groundwater on farms, though there is little evidence for it reaching deeper into local drinking water supplies. However, lagoon spills, such as from heavy rains in the wake of a hurricane, can lead to fish kills and algal blooms in local rivers. In the United States, of river across over 20 states were estimated to have been contaminated by manure leakage as of 2015. There is also evidence that evaporation from lagoons can cause nitrogen and phosphorus to spread through the air as dry particles then reach other water basins when they fall out through dry deposition. This process then also contributes to water eutrophication. On animal welfare Intensive pig production involves practices such as castration, earmarking, tattooing for litter identification, tail docking, which are often done without the use of anesthetic. Painful teeth clipping of piglets is also done to curtail cannibalism, behavioural instability and aggression, and tail biting, which are induced by the cramped environment. In English indoor farming, young pigs (less than 110kg in weight) are allowed to be kept with less than one square meter of space per pig. Pigs often begin life in a farrowing or gestation crate, which is a small pen with a central cage, designed to allow the piglets to feed from their mother while preventing her from attacking or crushing them. The crates are so small that the mother sows cannot turn around. While wild piglets remain with their mothers for around 12 to 14 weeks, farmed piglets are weaned and removed from their mothers at between two and five weeks old. Of the piglets born alive, 10% to 18% will not reach weaning age, instead succumbing to disease, starvation, dehydration, or accidental crushing by their mothers. Unusually small runt piglets are typically killed immediately by staff through blunt trauma to the head. Further, intensive farming involves sows giving birth to large litter sizes at an unnatural frequency, which increases the rate of stillborn piglets, and causes as many as 25%-50% of sows to die of prolapse. In culture Pigs, widespread in societies around the world since Neolithic times, have been used for many purposes in art, literature, and other expressions of human culture. In classical times, the Romans considered pork the finest of meats, enjoying sausages, and depicting them in their art. Across Europe, pigs have been celebrated in carnivals since the Middle Ages, becoming specially important in Medieval Germany in cities such as Nuremberg, and in Early Modern Italy in cities such as Bologna. Pigs, especially miniature breeds, are occasionally kept as pets. In literature, both for children and adults, pig characters appear in allegories, comic stories, and serious novels. In art, pigs have been represented in a wide range of media and styles from the earliest times in many cultures. Pig names are used in idioms and animal epithets, often derogatory, since pigs have long been linked with dirtiness and greed, while places such as Swindon are named for their association with swine. The eating of pork is forbidden in Islam and Judaism, but pigs are sacred in some other religions.	Biology and health sciences	Biology	null
2392005	https://en.wikipedia.org/wiki/Blum%20axioms	Blum axioms	In computational complexity theory the Blum axioms or Blum complexity axioms are axioms that specify desirable properties of complexity measures on the set of computable functions. The axioms were first defined by Manuel Blum in 1967. Importantly, Blum's speedup theorem and the Gap theorem hold for any complexity measure satisfying these axioms. The most well-known measures satisfying these axioms are those of time (i.e., running time) and space (i.e., memory usage). Definitions A Blum complexity measure is a pair with a numbering of the partial computable functions and a computable function which satisfies the following Blum axioms. We write for the i-th partial computable function under the Gödel numbering , and for the partial computable function . the domains of and are identical. the set is recursive. Examples is a complexity measure, if is either the time or the memory (or some suitable combination thereof) required for the computation coded by i. is not a complexity measure, since it fails the second axiom. Complexity classes For a total computable function complexity classes of computable functions can be defined as is the set of all computable functions with a complexity less than . is the set of all boolean-valued functions with a complexity less than . If we consider those functions as indicator functions on sets, can be thought of as a complexity class of sets.	Mathematics	Complexity theory	null
7713804	https://en.wikipedia.org/wiki/Rock-wallaby	Rock-wallaby	The rock-wallabies are the wallabies of the genus Petrogale. Taxonomy The genus was established in 1837 by John Edward Gray in a revision of material at the British Museum of Natural History. Gray nominated his earlier description of Kangurus pencillatus as the type species, now recognised in the combination Petrogale penicillata (brush-tailed rock-wallaby). The author separated the species from the defunct genus Kangurus, which he proposed to divide in his synopsis of the known macropod species. The following is a list of species, with common names, arranged by alliances of species groups: Genus Petrogale P. brachyotis species group Short-eared rock-wallaby, Petrogale brachyotis Monjon, Petrogale burbidgei Nabarlek, Petrogale concinna Eastern short-eared rock-wallaby, Petrogale wilkinsi P. xanthopus species group Proserpine rock-wallaby, Petrogale persephone Rothschild's rock-wallaby, Petrogale rothschildi Yellow-footed rock-wallaby, Petrogale xanthopus P. lateralis/penicillata species group Allied rock-wallaby, Petrogale assimilis Cape York rock-wallaby, Petrogale coenensis Godman's rock-wallaby, Petrogale godmani Herbert's rock-wallaby, Petrogale herberti Unadorned rock-wallaby, Petrogale inornata Black-flanked rock-wallaby, Petrogale lateralis Mareeba rock-wallaby, Petrogale mareeba Brush-tailed rock-wallaby, Petrogale penicillata Purple-necked rock-wallaby, Petrogale purpureicollis Mount Claro rock-wallaby, Petrogale sharmani Evolution and phylogenetics The species groups listed above have been confirmed by genetic analysis and their relationships have been well studied, especially in the brachyotis group. However, these studies also revealed that mitochondrial and nuclear DNA sequences resulted in different phylogenies, a phenomenon called cytonuclear discordance. Etymology From Latin petr- = rock + Greek galé = weasel. Description A genus with a high degree of speciation, driven in part by their fidelity to complex habitats that are phylogeographically isolated, Petrogale is the most diverse macropod genus, with workers identifying 19 species and further cryptic taxa in taxonomic revisions to 2014. The species occur in a weight range of 1–12 kilograms, relatively small to medium-sized marsupials. The medium-sized, often colourful and extremely agile rock-wallabies live where rocky, rugged and steep terrain can provide daytime refuge. Males are slightly larger than females, with a body length of up to 59 cm and a 70 cm long tail. Rock-wallabies are nocturnal and live a fortress-like existence spending their days in steep, rocky, complex terrain in some kind of shelter (a cave, an overhang or vegetation) and ranging out into surrounding terrain at night to feed. The greatest activity occurs three hours before sunrise and after sunset. Habitat Their reliance on refuges leads to the rock-wallabies living in small groups or colonies, with individuals having overlapping home ranges of about 15 hectares each. Within their colonies, they seem to be highly territorial with a male's territory overlapping one or a number of female territories. Even at night, the rock-wallabies do not move further than two kilometres from their home refuges. Generally, there are three categories of habitat that the different species of rock-wallaby seem to prefer: Loose piles of large boulders containing a maze of passageways Cliffs with many mid-level ledges and caves Isolated rock stacks, usually sheer sided and often girdled with fallen boulders Suitable habitat is limited and patchy and has led to varying degrees of isolation of colonies and a genetic differentiation specific to these colonies. The rock wallaby height is ranges from 60 cm to 70 cm. Species decline Their total numbers and range have been drastically reduced since European colonisation, with populations becoming extinct in the south. The ongoing extinction of colonies in recent times is of particular concern. In 1988 at Jenolan Caves in New South Wales, for example, a caged population of 80 rock-wallabies was released to boost what was thought to be an abundant local wild population. By 1992, the total population was down to about seven. The survivors were caught and enclosed in a fox and cat-proof enclosure, and the numbers in this captive population have since begun to increase. Scientists consider red foxes the major reason for the recent extinctions, along with competing herbivores, especially goats, sheep and rabbits, diseases such as toxoplasmosis and hydatidosis, habitat fragmentation and destruction, and a lower genetic health due to the increasing isolation of colonies. Recovery and conservation Habitat conservation and pest management addressing red foxes and goats appear to be the most urgent recovery actions to save the various species. The national recovery team with support from non-government organisations such as the Foundation for National Parks & Wildlife has implemented various programs ranging from land acquisition to captive breeding and awareness raising projects. Monitoring programs are implemented to register any changes in population sizes. Surveys and analysis establish the genetic diversity of populations. Red fox and goat eradication aid the survival of local populations, and captive breeding programs are used as an 'insurance policy' to build up rock-wallaby numbers to boost wild populations. In the case of the yellow-footed rock-wallaby, these strategies have prevented the extinction of the species in New South Wales.	Biology and health sciences	Diprotodontia	Animals
7715540	https://en.wikipedia.org/wiki/Chain	Chain	A chain is a serial assembly of connected pieces, called links, typically made of metal, with an overall character similar to that of a rope in that it is flexible and curved in compression but linear, rigid, and load-bearing in tension. A chain may consist of two or more links. Chains can be classified by their design, which can be dictated by their use: Those designed for lifting, such as when used with a hoist; for pulling; or for securing, such as with a bicycle lock, have links that are torus shaped, which make the chain flexible in two dimensions (the fixed third dimension being a chain's length). Small chains serving as jewellery are a mostly decorative analogue of such types. Those designed for transferring power in machines have links designed to mesh with the teeth of the sprockets of the machine, and are flexible in only one dimension. They are known as roller chains, though there are also non-roller chains such as block chains. Two distinct chains can be connected using a quick link, carabiner, shackle, or clevis. The load can be transferred from a chain to another object by a chain stopper. Uses for chains Uses for chains include: Decoration Belly chain, type of body jewelry worn around the waist Jewelry chain, many necklaces and bracelets are made out of small chains of gold and silver Chain of office, collar or heavy gold chain worn as insignia of office or a mark of fealty in medieval Europe and the United Kingdom Decorating clothing, some people wear wallets with chains connected to their belts, or pants decorated with chains Omega chain, a pseudo-chain where the 'links' are mounted on a backing rather than being interlinked Tie chain, used to hold in place a tie to the underlying shirt front Power transfer Bicycle chain, type of roller chain that transfers power from the pedals to the drive-wheel of a bicycle, thus propelling it. The chain is made up of a number of rigid links that are hinged together by pin joints to provide the flexibility needed to wrap around the bicycle's gears. Chain gun, type of machine gun that is driven by an external power source, sometimes connected by a chain, to actuate the mechanism rather than using recoil Chain pumps, type of water pump where a loop of chain inset discs is passed around then through a tube submerged in liquid Chainsaw, portable mechanical, motorized saw using a cutting chain to cut wood Timing chain, used to transfer rotational position from the crankshaft to the valve and ignition system on an internal combustion engine, typically with a 2:1 speed reduction. Security and restraint Ball and chain, a phrase that can refer to either the actual restraint device that was used to slow down prisoners, or a derogatory description of a person's significant other Belly chain (or waist chain), a physical restraint worn by prisoners, consisting of a chain around the prisoner's waist, to which the prisoner's hands are chained or cuffed Bicycle lock (or bicycle chain), lockable chain Chain boom, large chains used to exclude warships from harbors and rivers Chain link fencing, fencing that utilizes vertical wires that are bent in a zigzag fashion and linked to each other Chain mail, a type of armor consisting of small metal rings linked together in a pattern to form a mesh. Door chain, a type of security chain on a door that makes it possible to open a door from the inside while still making it difficult for someone outside to force their way inside Gang transport chain, a chain used to shackle two or more inmates together for transport or work outside the facility, forming a chain gang H-style restraints, a combination consisting of handcuffs on a belly chain with a connector chain running down to a set of leg irons Leg iron chains (fetters), an alternative to handcuffs Prisoner transport restraints, a combination which consists of a pair of handcuffs attached by a longer chain to a pair of leg irons On chain-linked handcuffs, the cuffs are held together by a short chain Traction, pulling and lifting Anchor cable, as used by ships and boats; in British nautical usage the component is a cable, the material is chain Chain slings Chain hoist, device used for lifting or lowering a load Chain boat, a type of river craft that used a steel chain laid along the riverbed for its propulsion Chain-linked lewis, a self-locking lifting device particularly for stone using a chain link as a pivot Curb chain, used on curb bits when riding a horse High-tensile chain (or transport chain), chain with a high tensile strength used for towing or securing loads Jack chain, a toothed chain used to move logs Lead shank (or stud chain), used on horses that are misbehaving Pull switch, an electrical switch operated by a ball chain Lavatory chain, the chain attached to the cistern of an old-fashioned W.C. in which the flushing power is obtained by a gravity feed from above-head height. Although most cisterns no longer work like that, the phrase "pull the chain" is still encountered to mean "flush the toilet". Rigid chain actuator, a type of chain that only bends in one direction, allowing it to operate under compression Snow chains, used to improve traction in snow Weapons Chain gun, type of machine gun that is driven by an external power source, sometimes connected by a chain, to actuate the mechanism rather than using recoil Chain shot, a type of ammunition for a cannon, used to inflict damage to the rigging of a sail vessel in naval warfare Chain weapon, a medieval weapon made of one or more weights attached to a handle with a chain Other uses Chains are a standard component of the deflection assembly of disc golf baskets. Chains can be used as a percussion instrument for special effects, such as in Arnold Schoenberg's Gurre-Lieder and Leoš Janáček's From the House of the Dead. Keychain, a small chain that connects a small item to a keyring Chain sinnet, a method of shortening a rope or other cable while in use or for storage Chain stitch, a sewing and embroidery technique Types of chain Ball chain, type of chain consisting of small sheet metal balls connected via short lengths of wire Calibrated chain, a type chain where the link lengths are within a given tolerance, so that it reliably engages with a windlass. Flat chain, form of chain used chiefly in agricultural machinery Ladder chain, a light wire chain used with sprockets for low torque power transmission Long link chain O-ring chain, a specialized type of roller chain Roller chain, the type of chain most commonly used for transmission of mechanical power on bicycles, motorcycles, and in industrial and agricultural machinery Self-lubricating chain, type of chain that uses a bush to continually lubricate the chain Silent chain, a type of chain in which the links engage the sprockets similarly to gear teeth Stud link chain, a type of chain with metal between the sides of each link, keeping the attached links in place. This helps prevent bunching when the chain is run out from a storage bin, as for use in anchoring ships. Short link chain, a chain where the gap between attached links is small relative to thickness. Connections Several methods are available to connect chain ends to each other or to other objects, and to apply a load to a chain away from the ends. These methods are usually specific to the type of chain, and must be of the correct size. Invention The metal link chain has been in use since at least 225 BC. Symbolism The prevalent modern symbolism is oppression, due to the use for a mechanical restriction of the liberty of a human or animal. Chains can also symbolize interconnectivity or interdependence. Unicode, in versions 6.x, contains the , which may show chain link(s). It may also denote a hyperlink. Gallery	Technology	Components_2	null
12382101	https://en.wikipedia.org/wiki/Drought%20tolerance	Drought tolerance	In botany, drought tolerance is the ability by which a plant maintains its biomass production during arid or drought conditions. Some plants are naturally adapted to dry conditions, surviving with protection mechanisms such as desiccation tolerance, detoxification, or repair of xylem embolism. Other plants, specifically crops like corn, wheat, and rice, have become increasingly tolerant to drought with new varieties created via genetic engineering. From an evolutionary perspective, the type of mycorrhizal associations formed in the roots of plants can determine how fast plants can adapt to drought. The plants behind drought tolerance are complex and involve many pathways which allows plants to respond to specific sets of conditions at any given time. Some of these interactions include stomatal conductance, carotenoid degradation and anthocyanin accumulation, the intervention of osmoprotectants (such as sucrose, glycine, and proline), ROS-scavenging enzymes. The molecular control of drought tolerance is also very complex and is influenced other factors such as environment and the developmental stage of the plant. This control consists mainly of transcriptional factors, such as dehydration-responsive element-binding protein (DREB), abscisic acid (ABA)-responsive element-binding factor (AREB), and NAM (no apical meristem). Physiology of drought tolerance Plants can be subjected to slowly developing water shortages (ie, taking days, weeks, or months), or they may face short-term deficits of water (ie, hours to days). In these situations, plants adapt by responding accordingly, minimizing water loss and maximizing water uptake. Plants are more susceptible to drought stress during the reproductive stages of growth, flowering and seed development. Therefore, the combination of short-term plus long-term responses allow for plants to produce a few viable seeds. Some examples of short-term and long-term physiological responses include: Short-term responses In the leaf: root-signal recognition, stomatal closure, decreased carbon assimilation In the stem: inhibition of growth, hydraulic changes, signal transport, assimilation of transport In the root: cell-drought signalling, osmotic adjustment Long-term responses In the above-ground portion of the plant: inhibition of shoot growth, reduced transpiration area, grain abortion, senescence, metabolic acclimation, osmotic adjustment, anthocyanin accumulation, carotenoid degradation, intervention of osmoprotectants, ROS-scavenging enzymes In the below-ground portion of the plant: turgor maintenance, sustained root growth, increased root/shoot, increased absorption area Regulatory network of drought tolerance In response to drought conditions, there is an alteration of gene expression, induced by or activated by transcription factors (TFs). These TFs bind to specific cis-elements to induce the expression of targeted stress-inducible genes, allowing for products to be transcribed that help with stress response and tolerance. Some of these include dehydration-responsive element-binding protein (DREB), ABA-responsive element-binding factor (AREB), no apical meristem (NAM), Arabidopsis transcription activation factor (ATAF), and cup-shaped cotyledon (CUC). Much of the molecular work to understand the regulation of drought tolerance has been done in Arabidopsis, helping elucidate the basic processes below. DREB TFs DREB1/CBF TFs DREB1A, DREB 1B, and DREB 1C are plant specific TFs which bind to drought responsive elements (DREs) in promoters responsive to drought, high salinity and low temperature in Arabidopsis. Overexpression of these genes enhance the tolerance of drought, high salinity, and low temperature in transgenic lines from Arabidopsis, rice, and tobacco. DEAR1/DREB and EAR motif protein 1 is a TF with an entirely different purpose nothing to do with drought stress. Tsutsui et al 2009 found Arabidopsis DEAR1 (At3g50260) to respond to pathogen infection, chitin, and oligomers of chitin. DREB2 TFs DREB proteins are involved in a variety of functions related to drought tolerance. For example, DREB proteins including DREB2A cooperate with AREB/ABF proteins in gene expression, specifically in the DREB2A gene under osmotic stress conditions. DREB2 also induces the expression of heat-related genes, such as heat shock protein. Overexpression of DREB2Aca enhances drought and heat stress tolerance levels in Arabidopsis. AREB/ABF TFs AREB/ABFs are ABA-responsive bZIP-type TFs which bind to ABA-responsive elements (ABREs) in stress-responsive promoters and activate gene expression. AREB1, AREB2, ABF3, and ABF1 have important roles in ABA signalling in the vegetative stage, as ABA controls the expression of genes associated with drought response and tolerance. The native form of AREB1 cannot target drought stress genes like RD29B in Arabidopsis, so modification is necessary for transcriptional activation. AREB/ABFs are positively regulated by SnRK2s, controlling the activity of target proteins via phosphorylation. This regulation also functions in the control of drought tolerance in the vegetative stage as well as the seed maturation and germination. Other TFs TFs such as NAC (composed of NAM, ATAF, and CUC), are also related to drought response in Arabidopsis and rice. Overexpression in the aforementioned plants improves stress and drought tolerance. They also may be related to root growth and senescence, two physiological traits related to drought tolerance. Natural drought tolerance adaptations Plants in naturally arid conditions retain large amounts of biomass due to drought tolerance and can be classified into 4 categories of adaptation: Drought-escaping plants: annuals that germinate and grow only during times of sufficient times of moisture to complete their life cycle. Drought-evading plants: non-succulent perennials which restrict their growth only to periods of moisture availability. Drought-enduring plants: also known as xerophytes, these evergreen shrubs have extensive root systems along with morphological and physiological adaptations which enable them to maintain growth even in times of extreme drought conditions. Drought-resisting plants: also known as succulent perennials, they have water stored in their leaves and stems for sparing uses. Structural adaptations Many adaptations for dry conditions are structural, including the following: Adaptations of the stomata to reduce water loss, such as reduced numbers, sunken pits, waxy surfaces.... Reduced number of leaves and their surface area. Water storage in succulent above-ground parts or water-filled tubers. Crassulacean acid metabolism (CAM metabolism) allows plants to get carbon dioxide at night and store malic acid during the day, allowing photosynthesis to take place with minimized water loss. Adaptations in the root system to increase water absorption. Trichomes (small hairs) on the leaves to absorb atmospheric water. Importance for agriculture With the frequency and severity of droughts increasing in recent years, damage to crops has become more serious, lowering the crop yield, growth, and production. However, research into the molecular pathways involving stress tolerance have revealed that overexpression of such genes can enhance drought tolerance, leading to projects focused on the development of transgenic crop varieties. Drought-tolerant plants which are developed through biotechnology enable farmers to protect their harvest and reduces losses in times of intense drought by using water more efficiently. Collaborations to improve drought tolerance in crop-variety plants International research projects to improve drought tolerance have been introduced, such as the Consultative Group on International Agricultural Research (CGIAR). One such project from CGIAR involves introducing genes such as DREB1 into lowland rice, upland rice, and wheat to evaluate drought tolerance in fields. This project aims to select at least 10 lines for agricultural use. Another similar project in collaboration with CGIAR, Embrapa, RIKEN, and the University of Tokyo have introduced AREB and DREB stress-tolerant genes into soybeans, finding several transgenic soybean lines with drought tolerance. Both projects have found improved grain yield and will be used to help develop future varieties that can be used commercially. Other examples of collaborations to improve drought tolerance in crop-variety plants include the International Center for Agricultural Research in Dry Areas (ICARDA) in Aleppo, Syria; the International Crops Research Institute for the Semi-Arid Tropics (ICRISAT) in Andhra Pradesh, India; the International Rice Research Institute (IRRI) in Los Baños, Philippines.; and the Heat and Drought Wheat Improvement Consortium (HeDWIC), a network that facilitates global coordination of wheat research to adapt to a future with more severe weather extremes. Impediments to the agricultural commercialization of drought tolerant plants The development of genetically modified crops includes multiple patents for genes and promoters, such as the marker genes in a vector, as well as transformation techniques. Therefore, freedom-to-operate (FTO) surveys should be implemented in collaborations for developing drought tolerant crops. Large amounts of money are also needed for the development of genetically modified groups. To bring a new genetically modified crop into the commercial market, it has been estimated to cost USD 136 million over 13 years. This poses a problem for development, as only a small number of companies can afford to develop drought-tolerant crops, and it is difficult for research institutions to sustain funding for this period of time. Therefore, a multinational framework with more collaboration among multiple disciples is needed to sustain projects of this size. Importance in horticulture Plant transformation has been used to develop multiple drought resistant crop varieties, but only limited varieties of ornamental plants. This significant lag in development is due to the fact that more transgenic ornamental plants are being developed for other reasons than drought tolerance. However, abiotic stress resistance is being explored in ornamental plants by Ornamental Biosciences. Transgenic Petunias, Poinsettias, New Guinea Impatiens, and Geraniums are being evaluated for frost, drought, and disease resistance. This will allow for a wider range of environments in which these plants can grow. Drought-tolerant plants This is a list of selected plant families, species and/or genus that tolerate drought: Acacia Adenium Aeonium Agapanthus Agave Aloe Angelonia Anredera cordifolia Araujia sericifera Arctotheca calendula Asphodelus Asparagus aethiopicus Banksia Begonia Bougainvillea Buddleja davidii Bulbine Cactus Calibrachoa Carpobrotus Cestrum parqui Cistus albidus Cosmos Crassula Curio Drosanthemum Dudleya Echeveria Echinacea Eriogonum Eucalyptus Euphorbia milii Gaillardia Gazania rigens Gonialoe Graptopetalum Grevillea Haberlea Haloxylon ammodendron Haworthia Helianthus Impatiens hawkeri Ipomoea cairica Lantana camara Melaleuca Mesembryanthemum Nandina domestica Nerium oleander Olea europaea Pelargonium Petunia Portulaca Portulacaria afra Ramonda serbica Ricinus communis Rosemary Salvia Santolina Sedum Senecio angulatus Senecio elegans Tetradenia riparia Thunbergia alata Vinca Vitis vinifera Yucca	Technology	Horticulture	null
5885549	https://en.wikipedia.org/wiki/Tubifex	Tubifex	Tubifex is a cosmopolitan genus of tubificid annelids that inhabits the sediments of lakes, rivers and occasionally sewer lines. At least 13 species of Tubifex have been identified, with the exact number not certain, as the species are not easily distinguishable from each other. Reproduction Tubifex worms are hermaphroditic: each individual has both male (testes) and female (ovaries) organs in the same animal. These minute reproductive organs are attached to the ventral side of the body wall in the celomic cavity. In mature specimens, the reproductive organs are clearly found on the ventral side of the body. Copulation and cocoon formation Although the Tubifex worms are hermaphrodites, the male and female organs become mature at different times; thus self-fertilization is avoided, and cross-fertilization is encouraged. Two mature Tubifex worms undergo copulation by joining ventral and anterior surfaces together with their anterior ends pointing opposite directions. Thus, the spermathecal opening of each worm is nearer to the male apertures of another worm. The penial setae of one worm penetrate into the tissues of other worm and thus the conjugants are held together. At this stage, the sperm of one worm is passed into the spermathecae of the other worm. After copulation, they separate and begin to produce egg cases containing eggs, called cocoons. The cocoon is formed around the clitellum as a soft, box-like structure into which the ova and the sperm are deposited. Soon, the Tubifex worm withdraws its body from the egg case by its backward wriggling movements. Culturing Tubifex Tubifex are raised commercially, mainly for fish food: the reddish Tubifex tubifex. Tubifex can be easily cultured on mass scale in containers with 50–75 mm thick pond mud at the bottom, blended with decaying vegetable matter and masses of bran and bread. Continuous, mild water flow is to be maintained in the container, with a suitable drainage system. After the arrangement of the system, the container is inoculated with Tubifex worms which can be obtained from nearby muddy canals or sewage canals. Within 15 days, clusters of worms develop and can be removed with mud in masses. When worms come to the surface due to lack of oxygen, they are collected and washed under brisk stream of water to remove residual mud attached to their bodies. Live food Tubifex worms are often used as a live food for fish, especially tropical fish and certain other freshwater species. They have been a popular food for the aquarium trade almost since its inception, and gathering them from open sewers for this purpose was quite common until recently. Most are now commercially obtained from the effluent of fish hatcheries, or from professional worm farms. Using these worms as a live food has come with certain problems over the years. When harvested from sewers, open bodies of water, and even from hatcheries, they may be infected with various diseases. This risk can be partially solved by keeping the worms under brisk running water until they have voided the contents of their digestive systems. However, the worms can still be vectors for whirling disease, which can affect salmonids. Additionally, they are very difficult for some fish to obtain in the wild, so certain fish, such as Rift Valley cichlids, will obsessively consume them until they make themselves sick. Additionally, while the worms have good-quality proteins, they also are very fattening, and are poor in certain important amino acids. Fish fed on them can grow rapidly, but may be less healthy and colorful than fish with more balanced diets. Lastly, in poorly cleaned aquaria, Tubifex can become established as a pest species, covering the bottom of the aquarium in a thick carpet which may be considered unsightly. Tubifex in sewers In 2009, a large blobby mass made of colonies of Tubifex was found to be living in the sewers of Raleigh, North Carolina. Revealed by a snake camera inspection of sewer piping under the Cameron Village shopping center, videos of the "creature" went viral on YouTube in 2009 under the name "Carolina poop monster". In 2013, an episode of the American television series Bar Rescue titled "Empty Bottles Full Cans" featured host Jon Taffer investigating the MT Bottle Bar in Tennessee, which had what the owners both claimed was a "natural spring" in the basement. The water there was in fact a natural underground spring without proper sewage pumping to filter it out from the basement and keep it from going stagnant; as a result, the water had turned to a black sludge-like consistency, and the episode featured multiple shock scenes where the camera focuses in on a large colony of Tubifex worms living there. The scene was later uploaded to YouTube under the title "Bar Rescue "Its natural spring water!" (MT Bottle)". Later in the episode, Taffer paid for a sewage system to irrigate and remove the excess water from the basement. Tubifex species The genus includes the following species: Tubifex blanchardi (Vejdovský, 1891) Tubifex harmani Loden, 1979 Tubifex costatus (Claparède, 1863) Tubifex ignotus (Stolc, 1886) Tubifex kryptus Bülow, 1957 Tubifex longipenis (Brinkhurst, 1965) Tubifex montanus Kowalewski, 1919 Tubifex nerthus Michaelsen, 1908 Tubifex newaensis (Michaelsen, 1903) Tubifex newfei Pickavance & Cook, 1971 Tubifex pescei (Dumnicka 1981) Tubifex pomoricus Timm, 1978 Tubifex pseudogaster (Dahl, 1960) Tubifex smirnowi Lastockin, 1927 Tubifex superiorensis Brinkhurst & Cook, 1971 Tubifex tubifex (O. F. Müller , 1774)	Biology and health sciences	Lophotrochozoa	Animals
207560	https://en.wikipedia.org/wiki/Photometry%20%28astronomy%29	Photometry (astronomy)	In astronomy, photometry, from Greek photo- ("light") and -metry ("measure"), is a technique used in astronomy that is concerned with measuring the flux or intensity of light radiated by astronomical objects. This light is measured through a telescope using a photometer, often made using electronic devices such as a CCD photometer or a photoelectric photometer that converts light into an electric current by the photoelectric effect. When calibrated against standard stars (or other light sources) of known intensity and colour, photometers can measure the brightness or apparent magnitude of celestial objects. The methods used to perform photometry depend on the wavelength region under study. At its most basic, photometry is conducted by gathering light and passing it through specialized photometric optical bandpass filters, and then capturing and recording the light energy with a photosensitive instrument. Standard sets of passbands (called a photometric system) are defined to allow accurate comparison of observations. A more advanced technique is spectrophotometry that is measured with a spectrophotometer and observes both the amount of radiation and its detailed spectral distribution. Photometry is also used in the observation of variable stars, by various techniques such as, differential photometry that simultaneously measures the brightness of a target object and nearby stars in the starfield or relative photometry by comparing the brightness of the target object to stars with known fixed magnitudes. Using multiple bandpass filters with relative photometry is termed absolute photometry. A plot of magnitude against time produces a light curve, yielding considerable information about the physical process causing the brightness changes. Precision photoelectric photometers can measure starlight around 0.001 magnitude. The technique of surface photometry can also be used with extended objects like planets, comets, nebulae or galaxies that measures the apparent magnitude in terms of magnitudes per square arcsecond. Knowing the area of the object and the average intensity of light across the astronomical object determines the surface brightness in terms of magnitudes per square arcsecond, while integrating the total light of the extended object can then calculate brightness in terms of its total magnitude, energy output or luminosity per unit surface area. Methods Astronomy was among the earliest applications of photometry. Modern photometers use specialised standard passband filters across the ultraviolet, visible, and infrared wavelengths of the electromagnetic spectrum. Any adopted set of filters with known light transmission properties is called a photometric system, and allows the establishment of particular properties about stars and other types of astronomical objects. Several important systems are regularly used, such as the UBV system (or the extended UBVRI system), near infrared JHK or the Strömgren uvbyβ system. Historically, photometry in the near-infrared through short-wavelength ultra-violet was done with a photoelectric photometer, an instrument that measured the light intensity of a single object by directing its light onto a photosensitive cell like a photomultiplier tube. These have largely been replaced with CCD cameras that can simultaneously image multiple objects, although photoelectric photometers are still used in special situations, such as where fine time resolution is required. Magnitudes and colour indices Modern photometric methods define magnitudes and colours of astronomical objects using electronic photometers viewed through standard coloured bandpass filters. This differs from other expressions of apparent visual magnitude observed by the human eye or obtained by photography: that usually appear in older astronomical texts and catalogues. Magnitudes measured by photometers in some commonplace photometric systems (UBV, UBVRI or JHK) are expressed with a capital letter, such as "V" (mV) or "B" (mB). Other magnitudes estimated by the human eye are expressed using lower case letters, such as "v", "b" or "p", etc. E.g. Visual magnitudes as mv, while photographic magnitudes are mph / mp or photovisual magnitudes mp or mpv. Hence, a 6th magnitude star might be stated as 6.0V, 6.0B, 6.0v or 6.0p. Because starlight is measured over a different range of wavelengths across the electromagnetic spectrum and are affected by different instrumental photometric sensitivities to light, they are not necessarily equivalent in numerical value. For example, apparent magnitude in the UBV system for the solar-like star 51 Pegasi is 5.46V, 6.16B or 6.39U, corresponding to magnitudes observed through each of the visual 'V', blue 'B' or ultraviolet 'U' filters. Magnitude differences between filters indicate colour differences and are related to temperature. Using B and V filters in the UBV system produces the B–V colour index. For 51 Pegasi, the B–V = 6.16 – 5.46 = +0.70, suggesting a yellow coloured star that agrees with its G2IV spectral type. Knowing the B–V results determines the star's surface temperature, finding an effective surface temperature of 5768±8 K. Another important application of colour indices is graphically plotting star's apparent magnitude against the B–V colour index. This forms the important relationships found between sets of stars in colour–magnitude diagrams, which for stars is the observed version of the Hertzsprung-Russell diagram. Typically photometric measurements of multiple objects obtained through two filters will show, for example in an open cluster, the comparative stellar evolution between the component stars or to determine the cluster's relative age. Due to the large number of different photometric systems adopted by astronomers, there are many expressions of magnitudes and their indices. Each of these newer photometric systems, excluding UBV, UBVRI or JHK systems, assigns an upper or lower case letter to the filter used. For example, magnitudes used by Gaia are 'G' (with the blue and red photometric filters, GBP and GRP) or the Strömgren photometric system having lower case letters of 'u', 'v', 'b', 'y', and two narrow and wide 'β' (Hydrogen-beta) filters. Some photometric systems also have certain advantages. For example, Strömgren photometry can be used to measure the effects of reddening and interstellar extinction. Strömgren allows calculation of parameters from the b and y filters (colour index of b − y) without the effects of reddening, as the indices m 1 and c 1. Applications There are many astronomical applications used with photometric systems. Photometric measurements can be combined with the inverse-square law to determine the luminosity of an object if its distance can be determined, or its distance if its luminosity is known. Other physical properties of an object, such as its temperature or chemical composition, may also be determined via broad or narrow-band spectrophotometry. Photometry is also used to study the light variations of objects such as variable stars, minor planets, active galactic nuclei and supernovae, or to detect transiting extrasolar planets. Measurements of these variations can be used, for example, to determine the orbital period and the radii of the members of an eclipsing binary star system, the rotation period of a minor planet or a star, or the total energy output of supernovae. CCD photometry A CCD (charge-coupled device) camera is essentially a grid of photometers, simultaneously measuring and recording the photons coming from all the sources in the field of view. Because each CCD image records the photometry of multiple objects at once, various forms of photometric extraction can be performed on the recorded data; typically relative, absolute, and differential. All three will require the extraction of the raw image magnitude of the target object, and a known comparison object. The observed signal from an object will typically cover many pixels according to the point spread function (PSF) of the system. This broadening is due to both the optics in the telescope and the astronomical seeing. When obtaining photometry from a point source, the flux is measured by summing all the light recorded from the object and subtracting the light due to the sky. The simplest technique, known as aperture photometry, consists of summing the pixel counts within an aperture centered on the object and subtracting the product of the nearby average sky count per pixel and the number of pixels within the aperture. This will result in the raw flux value of the target object. When doing photometry in a very crowded field, such as a globular cluster, where the profiles of stars overlap significantly, one must use de-blending techniques, such as PSF fitting to determine the individual flux values of the overlapping sources. Calibrations After determining the flux of an object in counts, the flux is normally converted into instrumental magnitude. Then, the measurement is calibrated in some way. Which calibrations are used will depend in part on what type of photometry is being done. Typically, observations are processed for relative or differential photometry. Relative photometry is the measurement of the apparent brightness of multiple objects relative to each other. Absolute photometry is the measurement of the apparent brightness of an object on a standard photometric system; these measurements can be compared with other absolute photometric measurements obtained with different telescopes or instruments. Differential photometry is the measurement of the difference in brightness of two objects. In most cases, differential photometry can be done with the highest precision, while absolute photometry is the most difficult to do with high precision. Also, accurate photometry is usually more difficult when the apparent brightness of the object is fainter. Absolute photometry To perform absolute photometry one must correct for differences between the effective passband through which an object is observed and the passband used to define the standard photometric system. This is often in addition to all of the other corrections discussed above. Typically this correction is done by observing the object(s) of interest through multiple filters and also observing a number of photometric standard stars. If the standard stars cannot be observed simultaneously with the target(s), this correction must be done under photometric conditions, when the sky is cloudless and the extinction is a simple function of the airmass. Relative photometry To perform relative photometry, one compares the instrument magnitude of the object to a known comparison object, and then corrects the measurements for spatial variations in the sensitivity of the instrument and the atmospheric extinction. This is often in addition to correcting for their temporal variations, particularly when the objects being compared are too far apart on the sky to be observed simultaneously. When doing the calibration from an image that contains both the target and comparison objects in close proximity, and using a photometric filter that matches the catalog magnitude of the comparison object most of the measurement variations decrease to null. Differential photometry Differential photometry is the simplest of the calibrations and most useful for time series observations. When using CCD photometry, both the target and comparison objects are observed at the same time, with the same filters, using the same instrument, and viewed through the same optical path. Most of the observational variables drop out and the differential magnitude is simply the difference between the instrument magnitude of the target object and the comparison object (∆Mag = C Mag – T Mag). This is very useful when plotting the change in magnitude over time of a target object, and is usually compiled into a light curve. Surface photometry For spatially extended objects such as galaxies, it is often of interest to measure the spatial distribution of brightness within the galaxy rather than simply measuring the galaxy's total brightness. An object's surface brightness is its brightness per unit solid angle as seen in projection on the sky, and measurement of surface brightness is known as surface photometry. A common application would be measurement of a galaxy's surface brightness profile, meaning its surface brightness as a function of distance from the galaxy's center. For small solid angles, a useful unit of solid angle is the square arcsecond, and surface brightness is often expressed in magnitudes per square arcsecond. The diameter of galaxies are often defined by the size of the 25th magnitude isophote in the blue B-band. Forced photometry In forced photometry, measurements are conducted at a specified location rather than for a specified object. It is "forced" in the sense that a measurement can be taken even if there is no object visible (in the spectral band of interest) in the location being observed. Forced photometry allows extracting a magnitude, or an upper limit for the magnitude, at a chosen sky location. Software A number of free computer programs are available for synthetic aperture photometry and PSF-fitting photometry. SExtractor and Aperture Photometry Tool are popular examples for aperture photometry. The former is geared towards reduction of large scale galaxy-survey data, and the latter has a graphical user interface (GUI) suitable for studying individual images. DAOPHOT is recognized as the best software for PSF-fitting photometry. Organizations There are a number of organizations, from professional to amateur, that gather and share photometric data and make it available on-line. Some sites gather the data primarily as a resource for other researchers (ex. AAVSO) and some solicit contributions of data for their own research (ex. CBA): American Association of Variable Star Observers (AAVSO). Astronomyonline.org Center for Backyard Astrophysics (CBA).	Physical sciences	Astrometry	null
207600	https://en.wikipedia.org/wiki/Hemiptera	Hemiptera	Hemiptera (; ) is an order of insects, commonly called true bugs, comprising over 80,000 species within groups such as the cicadas, aphids, planthoppers, leafhoppers, assassin bugs, bed bugs, and shield bugs. They range in size from to around , and share a common arrangement of piercing-sucking mouthparts. The name "true bugs" is often limited to the suborder Heteroptera. Entomologists reserve the term bug for Hemiptera or Heteroptera, which does not include other arthropods or insects of other orders such as ants, bees, beetles, or butterflies. In some varieties of English, all terrestrial arthropods (including non-insect arachnids and myriapods) also fall under the colloquial understanding of bug. Many insects with "bug" in their common name, especially in American English, belong to other orders; for example, the lovebug is a fly and the Maybug and ladybug are beetles. The term is occasionally extended to colloquial names for freshwater or marine crustaceans (e.g. Balmain bug, Moreton Bay bug, mudbug) and used by physicians and bacteriologists for disease-causing germs (e.g. superbugs). Most hemipterans feed on plants, using their sucking and piercing mouthparts to extract plant sap. Some are bloodsucking, or hematophagous, while others are predators that feed on other insects or small invertebrates. They live in a wide variety of habitats, generally terrestrial, though some are adapted to life in or on the surface of fresh water (e.g. pondskaters, water boatmen, giant water bugs). Hemipterans are hemimetabolous, with young nymphs that somewhat resemble adults. Many aphids are capable of parthenogenesis, producing young from unfertilised eggs; this helps them to reproduce extremely rapidly in favourable conditions. Humans have interacted with the Hemiptera for millennia. Some species, including many aphids, are significant agricultural pests, damaging crops by sucking the sap. Others harm humans more directly as vectors of serious viral diseases. The bed bug is a persistent parasite of humans, and some kissing bugs can transmit Chagas disease. Some species have been used for biological control of insect pests or of invasive plants. A few hemipterans, have been cultivated for the extraction of dyestuffs such as cochineal and carmine, and for shellac. Cicadas have been used as food, and have appeared in literature since the Iliad in Ancient Greece. Diversity Hemiptera is the largest order of hemimetabolous insects (not undergoing complete metamorphosis; though some examples such as male scale insects do undergo a form of complete metamorphosis ), containing over 95,000 named species. Other insect orders with more species are all holometabolous, meaning they have a pupal stage and undergo complete metamorphosis. The majority of species are terrestrial, including a number of important agricultural pests, but some are found in freshwater habitats. These include the water boatmen, backswimmers, pond skaters, and giant water bugs. Taxonomy and phylogeny Hemiptera belong to the insect superorder Paraneoptera, which includes lice (Psocodea), thrips (Thysanoptera), and the true bugs of Hemiptera. Within Paraneoptera, Hemiptera is most closely related to the sister clade Thysanoptera. The fossil record of hemipterans goes back to the Carboniferous (Moscovian). The oldest fossils are of the Archescytinidae from the Lower Permian and are thought to be basal to the Auchenorrhyncha. Fulgoromorpha and Cicadomorpha appear in the Upper Permian, as do Sternorrhyncha of the Psylloidea and Aleyrodoidea. Aphids and Coccoids appear in the Triassic. The Coleorrhyncha extend back to the Lower Jurassic. The Heteroptera first appeared in the Triassic. The present members of the order Hemiptera (sometimes referred to as Rhynchota) were historically placed into two orders, the so-called Homoptera and Heteroptera/Hemiptera, based on differences in wing structure and the position of the rostrum. "Homoptera" was established as paraphyletic group and an obsolete name. The order is now divided into four suborders, Heteroptera, Sternorrhyncha, Auchenorrhyncha, and Coleorrhyncha. The earlier work was based on nuclear DNA, but later phylogenetic analysis using mitochondrial DNA suggests that Homoptera may be monophyletic after all, a sister group of Heteroptera. The cause of the disparity in the analyses is suggested to be the long branch attraction effect in phylogenetic analysis, due to rapidly evolving DNA regions. The cladogram shows Hemiptera's placement within Paraneoptera, as well as how Hemiptera's four suborders are related. English names are given in parentheses where possible. Biology Mouthparts The defining feature of hemipterans is their "beak" in which the modified mandibles and maxillae form a "stylet" which is sheathed within a modified labium. The stylet is capable of piercing tissues and sucking liquids, while the labium supports it. The stylet contains a channel for the outward movement of saliva and another for the inward movement of liquid food. A salivary pump drives saliva into the prey; a cibarial pump extracts liquid from the prey. Both pumps are powered by substantial dilator muscles in the head. The beak is usually folded under the body when not in use. The diet is typically plant sap, but some hemipterans such as assassin bugs are predators. Both herbivorous and predatory hemipterans inject enzymes to begin digestion extra-orally (before the food is taken into the body). These enzymes include amylase to hydrolyse starch, polygalacturonase to weaken the tough cell walls of plants, and proteinases to break down proteins. Although the Hemiptera vary widely in their overall form, their mouthparts form a distinctive "rostrum". Other insect orders with mouthparts modified into anything like the rostrum and stylets of the Hemiptera include some Phthiraptera, but for other reasons they generally are easy to recognize as non-hemipteran. Similarly, the mouthparts of Siphonaptera, some Diptera and Thysanoptera superficially resemble the rostrum of the Hemiptera, but on closer inspection the differences are considerable. Aside from the mouthparts, various other insects can be confused with Hemiptera, but they all have biting mandibles and maxillae instead of the rostrum. Examples include cockroaches and psocids, both of which have longer, many-segmented antennae, and some beetles, but these have fully hardened forewings which do not overlap. Wing structure The forewings of Hemiptera are either entirely membranous, as in the Sternorrhyncha and Auchenorrhyncha, or partially hardened, as in most Heteroptera. The name "Hemiptera" is from the Greek (; "half") and (; "wing"), referring to the forewings of many heteropterans which are hardened near the base, but membranous at the ends. Wings modified in this manner are termed hemelytra (singular: hemelytron), by analogy with the completely hardened elytra of beetles, and occur only in the suborder Heteroptera. In all suborders, the hindwings – if present at all – are entirely membranous and usually shorter than the forewings. The forewings may be held "roofwise" over the body (typical of Sternorrhyncha and Auchenorrhyncha), or held flat on the back, with the ends overlapping (typical of Heteroptera). The antennae in Hemiptera typically consist of four or five segments, although they can still be quite long, and the tarsi of the legs have two or three segments. Sound production Many hemipterans can produce sound for communication. The "song" of male cicadas, the loudest of any insect, is produced by tymbal organs on the underside of the abdomen, and is used to attract mates. The tymbals are drumlike disks of cuticle, which are clicked in and out repeatedly, making a sound in the same way as popping the metal lid of a jam jar in and out. Stridulatory sounds are produced among the aquatic Corixidae and Notonectidae (backswimmers) using tibial combs rubbed across rostral ridges. Life cycle Hemipterans are hemimetabolous, meaning that they do not undergo metamorphosis, the complete change of form between a larval phase and an adult phase. Instead, their young are called nymphs, and resemble the adults to a greater or lesser degree. The nymphs moult several times as they grow, and each instar resembles the adult more than the previous one. Wing buds grow in later stage nymphs; the final transformation involves little more than the development of functional wings (if they are present at all) and functioning sexual organs, with no intervening pupal stage as in holometabolous insects. Parthenogenesis and viviparity Many aphids are parthenogenetic during part of the life cycle, such that females can produce unfertilized eggs, which are clones of their mother. All such young are females (thelytoky), so 100% of the population at these times can produce more offspring. Many species of aphid are also viviparous: the young are born live rather than laid as eggs. These adaptations enable aphids to reproduce extremely rapidly when conditions are suitable. Locomotion Hemipterans make use of a variety of modes of locomotion including swimming, skating on a water surface and jumping, as well as walking and flying like other insects. Swimming and skating Several families of Heteroptera are water bugs, adapted to an aquatic lifestyle, such as the water boatmen (Corixidae), water scorpions (Nepidae), and backswimmers (Notonectidae). They are mostly predatory, and have legs adapted as paddles to help the animal move through the water. The pondskaters or water striders (Gerridae) are also associated with water, but use the surface tension of standing water to keep them above the surface; they include the sea skaters in the genus Halobates, the only truly marine group of insects. Marangoni propulsion Marangoni effect propulsion exploits the change in surface tension when a soap-like surfactant is released on to a water surface, in the same way that a toy soap boat propels itself. Water bugs in the genus Microvelia (Veliidae) can travel at up to 17 cm/s, twice as fast as they can walk, by this means. Flight Flight is well developed in the Hemiptera although mostly used for short distance movement and dispersal. Wing development is sometimes related to environmental conditions. In some groups of Hemiptera, there are variations of winged, short-winged, and wingless forms within a single species. This kind of polymorphism tends to be helpful when habitats are temporary with more energy put into reproduction when food is available and into dispersal through flight when food becomes scarce. In aphids, both winged and wingless forms occur with winged forms produced in greater numbers when food resources are depleted. Aphids and whiteflies can sometimes be transported very long distances by atmospheric updrafts and high altitude winds. Wing-length polymorphism is notably rare in tree-living Hemiptera. Jumping Many Auchenorrhyncha including representatives of the cicadas, leafhoppers, treehoppers, planthoppers, and froghoppers are adapted for jumping (saltation). Treehoppers, for example, jump by rapidly depressing their hind legs. Before jumping, the hind legs are raised and the femora are pressed tightly into curved indentations in the coxae. Treehoppers can attain a take-off velocity of up to 2.7 metres per second and an acceleration of up to 250 g. The instantaneous power output is much greater than that of normal muscle, implying that energy is stored and released to catapult the insect into the air. Cicadas, which are much larger, extend their hind legs for a jump in under a millisecond, again implying elastic storage of energy for sudden release. Sedentary Instead of relying on any form of locomotion, most Sternorrhyncha females are sedentary or completely sessile, attached to their host plants by their thin feeding stylets which cannot be taken out of the plant quickly. Ecological roles Feeding modes Herbivores Most hemipterans are phytophagous, using their sucking and piercing mouthparts to feed on plant sap. These include cicadas, leafhoppers, treehoppers, planthoppers, froghoppers, aphids, whiteflies, scale insects, and some other groups. Some are monophages, being host specific and only found on one plant taxon, others are oligophages, feeding on a few plant groups, while others again are less discriminating polyphages and feed on many species of plant. The relationship between hemipterans and plants appears to be ancient, with piercing and sucking of plants evident in the Early Devonian period. Hemipterans can dramatically cut the mass of affected plants, especially in major outbreaks. They sometimes also change the mix of plants by predation on seeds or feeding on roots of certain species. Some sap-suckers move from one host to another at different times of year. Many aphids spend the winter as eggs on a woody host plant and the summer as parthenogenetically reproducing females on a herbaceous plant. Phloem sap, which has a higher concentration of sugars and nitrogen, is under positive pressure unlike the more dilute xylem sap. Most of the Sternorrhyncha and a number of Auchenorrhynchan groups feed on phloem. Phloem feeding is common in the Fulgoromorpha, most Cicadellidae and in the Heteroptera. The Typhlocybine Cicadellids specialize in feeding on non-vascular mesophyll tissue of leaves, which is more nutritious than the leaf epidermis. Most Heteroptera also feed on mesophyll tissue where they are more likely to encounter defensive secondary plant metabolites which often leads to the evolution of host specificity. Obligate xylem feeding is a special habit that is found in the Auchenorrhyncha among Cicadoidea, Cercopoidea and in Cicadelline Cicadellids. Some phloem feeders may take to xylem sap facultatively, especially when facing dehydration. Xylem feeders tend to be polyphagous; to overcome the negative pressure of xylem requires a special cibarial pump. Phloem feeding hemiptera typically have symbiotic micro-organisms in their gut that help to convert amino acids. Phloem feeders produce honeydew from their anus. A variety of organisms that feed on honeydew form symbiotic associations with phloem-feeders. Phloem sap is a sugary liquid low in amino acids, so insects have to process large quantities to meet their nutritional requirements. Xylem sap is even lower in amino acids and contains monosaccharides rather than sucrose, as well as organic acids and minerals. No digestion is required (except for the hydrolysis of sucrose) and 90% of the nutrients in the xylem sap can be utilised. Some phloem sap feeders selectively mix phloem and xylem sap to control the osmotic potential of the liquid consumed. A striking adaptation to a very dilute diet is found in many hemipterans: a filter chamber, a part of the gut looped back on itself as a countercurrent exchanger, which permits nutrients to be separated from excess water. The residue, mostly water with sugars and amino acids, is quickly excreted as sticky "honey dew", notably from aphids but also from other Auchenorrhycha and Sternorrhyncha. Some Sternorrhyncha including Psyllids and some aphids are gall formers. These sap-sucking hemipterans inject fluids containing plant hormones into the plant tissues inducing the production of tissue that covers to protects the insect and also act as sinks for nutrition that they feed on. The hackleberry gall psyllid for example, causes a woody gall on the leaf petioles of the hackleberry tree it infests, and the nymph of another psyllid produces a protective lerp out of hardened honeydew. Predators Most other hemipterans are predatory, feeding on other insects, or even small vertebrates. This is true of many aquatic species which are predatory, either as nymphs or adults. The predatory shield bug for example stabs caterpillars with its beak and sucks out the body fluids. The saliva of predatory heteropterans contains digestive enzymes such as proteinase and phospholipase, and in some species also amylase. The mouthparts of these insects are adapted for predation. There are toothed stylets on the mandibles able to cut into and abrade tissues of their prey. There are further stylets on the maxillae, adapted as tubular canals to inject saliva and to extract the pre-digested and liquified contents of the prey. Haematophagic ectoparasites A few hemipterans are haematophagic ectoparasites), feeding on the blood of larger animals. These include bedbugs and the triatomine kissing bugs of the assassin bug family Reduviidae, which can transmit the dangerous Chagas disease. The first known hemipteran to feed in this way on vertebrates was the extinct assassin bug Triatoma dominicana found fossilized in amber and dating back about twenty million years. Faecal pellets fossilised beside it show that it transmitted a disease-causing Trypanosoma and the amber included hairs of the likely host, a bat. As symbionts Some species of ant protect and farm aphids (Sternorrhyncha) and other sap-sucking hemipterans, gathering and eating the honeydew that these hemipterans secrete. The relationship is mutualistic, as both ant and aphid benefit. Ants such as the yellow anthill ant, Lasius flavus, breed aphids of at least four species, Geoica utricularia, Tetraneura ulmi, Forda marginata and Forda formicaria, taking eggs with them when they found a new colony; in return, these aphids are obligately associated with the ant, breeding mainly or wholly asexually inside anthills. Ants may also protect the plant bugs from their natural enemies, removing the eggs of predatory beetles and preventing access by parasitic wasps. Some leafhoppers (Auchenorrhyncha) are similarly "milked" by ants. In the Corcovado rain forest of Costa Rica, wasps compete with ants to protect and milk leafhoppers; the leafhoppers preferentially give more honeydew, more often, to the wasps, which are larger and may offer better protection. As prey: defences against predators and parasites Hemiptera form prey to predators including vertebrates, such as birds, and other invertebrates such as ladybirds. In response, hemipterans have evolved antipredator adaptations. Ranatra may feign death (thanatosis). Others such as Carpocoris purpureipennis secrete toxic fluids to ward off arthropod predators; some Pentatomidae such as Dolycoris are able to direct these fluids at an attacker. Toxic cardenolide compounds are accumulated by the heteropteran Oncopeltus fasciatus when it consumes milkweeds, while the coreid stinkbug Amorbus rubiginosus acquires 2-hexenal from its food plant, Eucalyptus. Some long-legged bugs mimic twigs, rocking to and fro to simulate the motion of a plant part in the wind. The nymph of the Masked hunter bug camouflages itself with sand grains, using its hind legs and tarsal fan to form a double layer of grains, coarser on the outside. The Amazon rain forest cicada Hemisciera maculipennis displays bright red deimatic flash coloration on its hindwings when threatened; the sudden contrast helps to startle predators, giving the cicada time to escape. The coloured patch on the hindwing is concealed at rest by an olive green patch of the same size on the forewing, enabling the insect to switch rapidly from cryptic to deimatic behaviour. Some hemipterans such as firebugs have bold aposematic warning coloration, often red and black, which appear to deter passerine birds. Many hemipterans including aphids, scale insects and especially the planthoppers secrete wax to protect themselves from threats such as fungi, parasitoidal insects and predators, as well as abiotic factors like desiccation. Hard waxy coverings are especially important in the sedentary Sternorrhyncha such as scale insects, which have no means of escaping from predators; other Sternorrhyncha evade detection and attack by creating and living inside plant galls. Nymphal Cicadoidea and Cercopoidea have glands attached to the Malpighian tubules in their proximal segment that produce mucopolysaccharides, which form the froth around spittlebugs, offering a measure of protection. Parental care is found in many species of Hemiptera especially in members of the Membracidae and numerous Heteroptera. In many species of shield bug, females stand guard over their egg clusters to protect them from egg parasitoids and predators. In the aquatic Belostomatidae, females lay their eggs on the back of the male which guards the eggs. Protection provided by ants is common in the Auchenorrhyncha. Interaction with humans As pests Although many species of Hemiptera are significant pests of crops and garden plants, including many species of aphid and scale insects, other species are harmless. The damage done is often not so much the deprivation of the plant of its sap, but the fact that they transmit serious viral diseases between plants. They often produce copious amounts of honeydew which encourages the growth of sooty mould. Significant pests include the cottony cushion scale, a pest of citrus fruit trees, the green peach aphid and other aphids which attack crops worldwide and transmit diseases, and jumping plant lice which can be plant-specific and transmit diseases, as with the Asian citrus psyllid which transmits citrus greening disease. For pest control Members of the families Reduviidae, Phymatidae and Nabidae are obligate predators. Some predatory species are used in biological pest control; these include various nabids, and even some members of families that are primarily phytophagous, such as the genus Geocoris in the family Lygaeidae. Other hemipterans are omnivores, alternating between a plant-based and an animal-based diet. For example, Dicyphus hesperus is used to control whitefly on tomatoes but also sucks sap, and if deprived of plant tissues will die even if in the presence of whiteflies. The spined soldier bug (Podisus maculiventris) sucks body fluids from several pests including the larvae of the Colorado beetle and the Mexican bean beetle. Insect products Other hemipterans have positive uses for humans, such as in the production of the dyestuff carmine (cochineal). The FDA has created guidelines for how to declare when it has been added to a product. The scale insect Dactylopius coccus produces the brilliant red-coloured carminic acid to deter predators. Up to 100,000 scale insects need to be collected and processed to make a kilogram (2.2 lbs) of cochineal dye. A similar number of lac bugs are needed to make a kilogram of shellac, a brush-on colourant and wood finish. Additional uses of this traditional product include the waxing of citrus fruits to extend their shelf-life, and the coating of pills to moisture-proof them, provide slow-release or mask the taste of bitter ingredients. As human parasites and disease vectors Chagas disease is a modern-day tropical disease caused by Trypanosoma cruzi and transmitted by kissing bugs, so-called because they suck human blood from around the lips while a person sleeps. The bed bug, Cimex lectularius, is an external parasite of humans. It lives in bedding and is mainly active at night, feeding on human blood, generally without being noticed. Bed bugs mate by traumatic insemination; the male pierces the female's abdomen and injects his sperm into a secondary genital structure, the spermalege. The sperm travel in the female's blood (haemolymph) to sperm storage structures (seminal conceptacles); they are released from there to fertilise her eggs inside her ovaries. As food Some larger hemipterans such as cicadas are used as food in Asian countries such as China, and they are much esteemed in Malawi and other African countries. Insects have a high protein content and good food conversion ratios, but most hemipterans are too small to be a useful component of the human diet. At least nine species of Hemiptera are eaten worldwide. In art and literature Cicadas have featured in literature since the time of Homer's Iliad, and as motifs in decorative art from the Chinese Shang dynasty (1766–1122 B.C.). They are described by Aristotle in his History of Animals and by Pliny the Elder in his Natural History; their mechanism of sound production is mentioned by Hesiod in his poem Works and Days "when the Skolymus flowers, and the tuneful Tettix sitting on his tree in the weary summer season pours forth from under his wings his shrill song". In mythology and folklore Among the bugs, cicadas in particular have been used as money, in folk medicine, to forecast the weather, to provide song (in China), and in folklore and myths around the world. Threats Large-scale cultivation of the oil palm Elaeis guineensis in the Amazon basin damages freshwater habitats and reduces the diversity of aquatic and semi-aquatic Heteroptera. Climate change may be affecting the global migration of hemipterans including the potato leafhopper, Empoasca fabae. Warming is correlated with the severity of potato leafhopper infestation, so increased warming may worsen infestations in future.	Biology and health sciences	Hemiptera (true bugs)	null
207601	https://en.wikipedia.org/wiki/Adsorption	Adsorption	Adsorption is the adhesion of atoms, ions or molecules from a gas, liquid or dissolved solid to a surface. This process creates a film of the adsorbate on the surface of the adsorbent. This process differs from absorption, in which a fluid (the absorbate) is dissolved by or permeates a liquid or solid (the absorbent). While adsorption does often precede absorption, which involves the transfer of the absorbate into the volume of the absorbent material, alternatively, adsorption is distinctly a surface phenomenon, wherein the adsorbate does not penetrate through the material surface and into the bulk of the adsorbent. The term sorption encompasses both adsorption and absorption, and desorption is the reverse of sorption. Like surface tension, adsorption is a consequence of surface energy. In a bulk material, all the bonding requirements (be they ionic, covalent or metallic) of the constituent atoms of the material are fulfilled by other atoms in the material. However, atoms on the surface of the adsorbent are not wholly surrounded by other adsorbent atoms and therefore can attract adsorbates. The exact nature of the bonding depends on the details of the species involved, but the adsorption process is generally classified as physisorption (characteristic of weak van der Waals forces) or chemisorption (characteristic of covalent bonding). It may also occur due to electrostatic attraction. The nature of the adsorption can affect the structure of the adsorbed species. For example, polymer physisorption from solution can result in squashed structures on a surface. Adsorption is present in many natural, physical, biological and chemical systems and is widely used in industrial applications such as heterogeneous catalysts, activated charcoal, capturing and using waste heat to provide cold water for air conditioning and other process requirements (adsorption chillers), synthetic resins, increasing storage capacity of carbide-derived carbons and water purification. Adsorption, ion exchange and chromatography are sorption processes in which certain adsorbates are selectively transferred from the fluid phase to the surface of insoluble, rigid particles suspended in a vessel or packed in a column. Pharmaceutical industry applications, which use adsorption as a means to prolong neurological exposure to specific drugs or parts thereof, are lesser known. The word "adsorption" was coined in 1881 by German physicist Heinrich Kayser (1853–1940). Isotherms The adsorption of gases and solutes is usually described through isotherms, that is, the amount of adsorbate on the adsorbent as a function of its pressure (if gas) or concentration (for liquid phase solutes) at constant temperature. The quantity adsorbed is nearly always normalized by the mass of the adsorbent to allow comparison of different materials. To date, 15 different isotherm models have been developed. Freundlich The first mathematical fit to an isotherm was published by Freundlich and Kuster (1906) and is a purely empirical formula for gaseous adsorbates: where is the mass of adsorbate adsorbed, is the mass of the adsorbent, is the pressure of adsorbate (this can be changed to concentration if investigating solution rather than gas), and and are empirical constants for each adsorbent–adsorbate pair at a given temperature. The function is not adequate at very high pressure because in reality has an asymptotic maximum as pressure increases without bound. As the temperature increases, the constants and change to reflect the empirical observation that the quantity adsorbed rises more slowly and higher pressures are required to saturate the surface. Langmuir Irving Langmuir was the first to derive a scientifically based adsorption isotherm in 1918. The model applies to gases adsorbed on solid surfaces. It is a semi-empirical isotherm with a kinetic basis and was derived based on statistical thermodynamics. It is the most common isotherm equation to use due to its simplicity and its ability to fit a variety of adsorption data. It is based on four assumptions: All of the adsorption sites are equivalent, and each site can only accommodate one molecule. The surface is energetically homogeneous, and adsorbed molecules do not interact. There are no phase transitions. At the maximum adsorption, only a monolayer is formed. Adsorption only occurs on localized sites on the surface, not with other adsorbates. These four assumptions are seldom all true: there are always imperfections on the surface, adsorbed molecules are not necessarily inert, and the mechanism is clearly not the same for the first molecules to adsorb to a surface as for the last. The fourth condition is the most troublesome, as frequently more molecules will adsorb to the monolayer; this problem is addressed by the BET isotherm for relatively flat (non-microporous) surfaces. The Langmuir isotherm is nonetheless the first choice for most models of adsorption and has many applications in surface kinetics (usually called Langmuir–Hinshelwood kinetics) and thermodynamics. Langmuir suggested that adsorption takes place through this mechanism: , where A is a gas molecule, and S is an adsorption site. The direct and inverse rate constants are k and k−1. If we define surface coverage, , as the fraction of the adsorption sites occupied, in the equilibrium we have: or where is the partial pressure of the gas or the molar concentration of the solution. For very low pressures , and for high pressures . The value of is difficult to measure experimentally; usually, the adsorbate is a gas and the quantity adsorbed is given in moles, grams, or gas volumes at standard temperature and pressure (STP) per gram of adsorbent. If we call vmon the STP volume of adsorbate required to form a monolayer on the adsorbent (per gram of adsorbent), then , and we obtain an expression for a straight line: Through its slope and y intercept we can obtain vmon and K, which are constants for each adsorbent–adsorbate pair at a given temperature. vmon is related to the number of adsorption sites through the ideal gas law. If we assume that the number of sites is just the whole area of the solid divided into the cross section of the adsorbate molecules, we can easily calculate the surface area of the adsorbent. The surface area of an adsorbent depends on its structure: the more pores it has, the greater the area, which has a big influence on reactions on surfaces. If more than one gas adsorbs on the surface, we define as the fraction of empty sites, and we have: Also, we can define as the fraction of the sites occupied by the j-th gas: where i is each one of the gases that adsorb. Note: 1) To choose between the Langmuir and Freundlich equations, the enthalpies of adsorption must be investigated. While the Langmuir model assumes that the energy of adsorption remains constant with surface occupancy, the Freundlich equation is derived with the assumption that the heat of adsorption continually decrease as the binding sites are occupied. The choice of the model based on best fitting of the data is a common misconception. 2) The use of the linearized form of the Langmuir model is no longer common practice. Advances in computational power allowed for nonlinear regression to be performed quickly and with higher confidence since no data transformation is required. BET Often molecules do form multilayers, that is, some are adsorbed on already adsorbed molecules, and the Langmuir isotherm is not valid. In 1938 Stephen Brunauer, Paul Emmett, and Edward Teller developed a model isotherm that takes that possibility into account. Their theory is called BET theory, after the initials in their last names. They modified Langmuir's mechanism as follows: A(g) + S ⇌ AS, A(g) + AS ⇌ A2S, A(g) + A2S ⇌ A3S and so on. The derivation of the formula is more complicated than Langmuir's (see links for complete derivation). We obtain: where x is the pressure divided by the vapor pressure for the adsorbate at that temperature (usually denoted ), v is the STP volume of adsorbed adsorbate, vmon is the STP volume of the amount of adsorbate required to form a monolayer, and c is the equilibrium constant K we used in Langmuir isotherm multiplied by the vapor pressure of the adsorbate. The key assumption used in deriving the BET equation that the successive heats of adsorption for all layers except the first are equal to the heat of condensation of the adsorbate. The Langmuir isotherm is usually better for chemisorption, and the BET isotherm works better for physisorption for non-microporous surfaces. Kisliuk In other instances, molecular interactions between gas molecules previously adsorbed on a solid surface form significant interactions with gas molecules in the gaseous phases. Hence, adsorption of gas molecules to the surface is more likely to occur around gas molecules that are already present on the solid surface, rendering the Langmuir adsorption isotherm ineffective for the purposes of modelling. This effect was studied in a system where nitrogen was the adsorbate and tungsten was the adsorbent by Paul Kisliuk (1922–2008) in 1957. To compensate for the increased probability of adsorption occurring around molecules present on the substrate surface, Kisliuk developed the precursor state theory, whereby molecules would enter a precursor state at the interface between the solid adsorbent and adsorbate in the gaseous phase. From here, adsorbate molecules would either adsorb to the adsorbent or desorb into the gaseous phase. The probability of adsorption occurring from the precursor state is dependent on the adsorbate's proximity to other adsorbate molecules that have already been adsorbed. If the adsorbate molecule in the precursor state is in close proximity to an adsorbate molecule that has already formed on the surface, it has a sticking probability reflected by the size of the SE constant and will either be adsorbed from the precursor state at a rate of kEC or will desorb into the gaseous phase at a rate of kES. If an adsorbate molecule enters the precursor state at a location that is remote from any other previously adsorbed adsorbate molecules, the sticking probability is reflected by the size of the SD constant. These factors were included as part of a single constant termed a "sticking coefficient", kE, described below: As SD is dictated by factors that are taken into account by the Langmuir model, SD can be assumed to be the adsorption rate constant. However, the rate constant for the Kisliuk model (R’) is different from that of the Langmuir model, as R’ is used to represent the impact of diffusion on monolayer formation and is proportional to the square root of the system's diffusion coefficient. The Kisliuk adsorption isotherm is written as follows, where θ(t) is fractional coverage of the adsorbent with adsorbate, and t is immersion time: Solving for θ(t) yields: Adsorption enthalpy Adsorption constants are equilibrium constants, therefore they obey the Van 't Hoff equation: As can be seen in the formula, the variation of K must be isosteric, that is, at constant coverage. If we start from the BET isotherm and assume that the entropy change is the same for liquefaction and adsorption, we obtain that is to say, adsorption is more exothermic than liquefaction. Single-molecule explanation The adsorption of ensemble molecules on a surface or interface can be divided into two processes: adsorption and desorption. If the adsorption rate wins the desorption rate, the molecules will accumulate over time giving the adsorption curve over time. If the desorption rate is larger, the number of molecules on the surface will decrease over time. The adsorption rate is dependent on the temperature, the diffusion rate of the solute (related to mean free path for pure gas), and the energy barrier between the molecule and the surface. The diffusion and key elements of the adsorption rate can be calculated using Fick's laws of diffusion and Einstein relation (kinetic theory). Under ideal conditions, when there is no energy barrier and all molecules that diffuse and collide with the surface get adsorbed, the number of molecules adsorbed at a surface of area on an infinite area surface can be directly integrated from Fick's second law differential equation to be: where is the surface area (unit m2), is the number concentration of the molecule in the bulk solution (unit #/m3), is the diffusion constant (unit m2/s), and is time (unit s). Further simulations and analysis of this equation show that the square root dependence on the time is originated from the decrease of the concentrations near the surface under ideal adsorption conditions. Also, this equation only works for the beginning of the adsorption when a well-behaved concentration gradient forms near the surface. Correction on the reduction of the adsorption area and slowing down of the concentration gradient evolution have to be considered over a longer time. Under real experimental conditions, the flow and the small adsorption area always make the adsorption rate faster than what this equation predicted, and the energy barrier will either accelerate this rate by surface attraction or slow it down by surface repulsion. Thus, the prediction from this equation is often a few to several orders of magnitude away from the experimental results. Under special cases, such as a very small adsorption area on a large surface, and under chemical equilibrium when there is no concentration gradience near the surface, this equation becomes useful to predict the adsorption rate with debatable special care to determine a specific value of in a particular measurement. The desorption of a molecule from the surface depends on the binding energy of the molecule to the surface and the temperature. The typical overall adsorption rate is thus often a combined result of the adsorption and desorption. Quantum mechanical – thermodynamic modelling for surface area and porosity Since 1980 two theories were worked on to explain adsorption and obtain equations that work. These two are referred to as the chi hypothesis, the quantum mechanical derivation, and excess surface work (ESW). Both these theories yield the same equation for flat surfaces: where U is the unit step function. The definitions of the other symbols is as follows: where "ads" stands for "adsorbed", "m" stands for "monolayer equivalence" and "vap" is reference to the vapor pressure of the liquid adsorptive at the same temperature as the solid sample. The unit function creates the definition of the molar energy of adsorption for the first adsorbed molecule by: The plot of adsorbed versus is referred to as the chi plot. For flat surfaces, the slope of the chi plot yields the surface area. Empirically, this plot was noticed as being a very good fit to the isotherm by Michael Polanyi and also by Jan Hendrik de Boer and Cornelis Zwikker but not pursued. This was due to criticism in the former case by Albert Einstein and in the latter case by Brunauer. This flat surface equation may be used as a "standard curve" in the normal tradition of comparison curves, with the exception that the porous sample's early portion of the plot of versus acts as a self-standard. Ultramicroporous, microporous and mesoporous conditions may be analyzed using this technique. Typical standard deviations for full isotherm fits including porous samples are less than 2%. Notice that in this description of physical adsorption, the entropy of adsorption is consistent with the Dubinin thermodynamic criterion, that is the entropy of adsorption from the liquid state to the adsorbed state is approximately zero. Adsorbents Characteristics and general requirements Adsorbents are used usually in the form of spherical pellets, rods, moldings, or monoliths with a hydrodynamic radius between 0.25 and 5 mm. They must have high abrasion resistance, high thermal stability and small pore diameters, which results in higher exposed surface area and hence high capacity for adsorption. The adsorbents must also have a distinct pore structure that enables fast transport of the gaseous vapors. Most industrial adsorbents fall into one of three classes: Oxygen-containing compounds – Are typically hydrophilic and polar, including materials such as silica gel, limestone (calcium carbonate) and zeolites. Carbon-based compounds – Are typically hydrophobic and non-polar, including materials such as activated carbon and graphite. Polymer-based compounds – Are polar or non-polar, depending on the functional groups in the polymer matrix. Silica gel Silica gel is a chemically inert, non-toxic, polar and dimensionally stable (< ) amorphous form of SiO2. It is prepared by the reaction between sodium silicate and acetic acid, which is followed by a series of after-treatment processes such as aging, pickling, etc. These after-treatment methods results in various pore size distributions. Silica is used for drying of process air (e.g. oxygen, natural gas) and adsorption of heavy (polar) hydrocarbons from natural gas. Zeolites Zeolites are natural or synthetic crystalline aluminosilicates, which have a repeating pore network and release water at high temperature. Zeolites are polar in nature. They are manufactured by hydrothermal synthesis of sodium aluminosilicate or another silica source in an autoclave followed by ion exchange with certain cations (Na+, Li+, Ca2+, K+, NH4+). The channel diameter of zeolite cages usually ranges from 2 to 9 Å. The ion exchange process is followed by drying of the crystals, which can be pelletized with a binder to form macroporous pellets. Zeolites are applied in drying of process air, CO2 removal from natural gas, CO removal from reforming gas, air separation, catalytic cracking, and catalytic synthesis and reforming. Non-polar (siliceous) zeolites are synthesized from aluminum-free silica sources or by dealumination of aluminum-containing zeolites. The dealumination process is done by treating the zeolite with steam at elevated temperatures, typically greater than . This high temperature heat treatment breaks the aluminum-oxygen bonds and the aluminum atom is expelled from the zeolite framework. Activated carbon The term "adsorption" itself was coined by Heinrich Kayser in 1881 in the context of uptake of gases by carbons. Activated carbon is a highly porous, amorphous solid consisting of microcrystallites with a graphite lattice, usually prepared in small pellets or a powder. It is non-polar and cheap. One of its main drawbacks is that it reacts with oxygen at moderate temperatures (over 300 °C). Activated carbon can be manufactured from carbonaceous material, including coal (bituminous, subbituminous, and lignite), peat, wood, or nutshells (e.g., coconut). The manufacturing process consists of two phases, carbonization and activation. The carbonization process includes drying and then heating to separate by-products, including tars and other hydrocarbons from the raw material, as well as to drive off any gases generated. The process is completed by heating the material over in an oxygen-free atmosphere that cannot support combustion. The carbonized particles are then "activated" by exposing them to an oxidizing agent, usually steam or carbon dioxide at high temperature. This agent burns off the pore blocking structures created during the carbonization phase and so, they develop a porous, three-dimensional graphite lattice structure. The size of the pores developed during activation is a function of the time that they spend in this stage. Longer exposure times result in larger pore sizes. The most popular aqueous phase carbons are bituminous based because of their hardness, abrasion resistance, pore size distribution, and low cost, but their effectiveness needs to be tested in each application to determine the optimal product. Activated carbon is used for adsorption of organic substances and non-polar adsorbates and it is also usually used for waste gas (and waste water) treatment. It is the most widely used adsorbent since most of its chemical (e.g. surface groups) and physical properties (e.g. pore size distribution and surface area) can be tuned according to what is needed. Its usefulness also derives from its large micropore (and sometimes mesopore) volume and the resulting high surface area. Recent research works reported activated carbon as an effective agent to adsorb cationic species of toxic metals from multi-pollutant systems and also proposed possible adsorption mechanisms with supporting evidences. Water adsorption The adsorption of water at surfaces is of broad importance in chemical engineering, materials science and catalysis. Also termed surface hydration, the presence of physically or chemically adsorbed water at the surfaces of solids plays an important role in governing interface properties, chemical reaction pathways and catalytic performance in a wide range of systems. In the case of physically adsorbed water, surface hydration can be eliminated simply through drying at conditions of temperature and pressure allowing full vaporization of water. For chemically adsorbed water, hydration may be in the form of either dissociative adsorption, where H2O molecules are dissociated into surface adsorbed -H and -OH, or molecular adsorption (associative adsorption) where individual water molecules remain intact Adsorption solar heating and storage The low cost ($200/ton) and high cycle rate (2,000 ×) of synthetic zeolites such as Linde 13X with water adsorbate has garnered much academic and commercial interest recently for use for thermal energy storage (TES), specifically of low-grade solar and waste heat. Several pilot projects have been funded in the EU from 2000 to the present (2020). The basic concept is to store solar thermal energy as chemical latent energy in the zeolite. Typically, hot dry air from flat plate solar collectors is made to flow through a bed of zeolite such that any water adsorbate present is driven off. Storage can be diurnal, weekly, monthly, or even seasonal depending on the volume of the zeolite and the area of the solar thermal panels. When heat is called for during the night, or sunless hours, or winter, humidified air flows through the zeolite. As the humidity is adsorbed by the zeolite, heat is released to the air and subsequently to the building space. This form of TES, with specific use of zeolites, was first taught by John Guerra in 1978. Carbon capture and storage Typical adsorbents proposed for carbon capture and storage are zeolites and MOFs. The customization of adsorbents makes them a potentially attractive alternative to absorption. Because adsorbents can be regenerated by temperature or pressure swing, this step can be less energy intensive than absorption regeneration methods. Major problems that are present with adsorption cost in carbon capture are: regenerating the adsorbent, mass ratio, solvent/MOF, cost of adsorbent, production of the adsorbent, lifetime of adsorbent. In sorption enhanced water gas shift (SEWGS) technology a pre-combustion carbon capture process, based on solid adsorption, is combined with the water gas shift reaction (WGS) in order to produce a high pressure hydrogen stream. The CO2 stream produced can be stored or used for other industrial processes. Protein and surfactant adsorption Protein adsorption is a process that has a fundamental role in the field of biomaterials. Indeed, biomaterial surfaces in contact with biological media, such as blood or serum, are immediately coated by proteins. Therefore, living cells do not interact directly with the biomaterial surface, but with the adsorbed proteins layer. This protein layer mediates the interaction between biomaterials and cells, translating biomaterial physical and chemical properties into a "biological language". In fact, cell membrane receptors bind to protein layer bioactive sites and these receptor-protein binding events are transduced, through the cell membrane, in a manner that stimulates specific intracellular processes that then determine cell adhesion, shape, growth and differentiation. Protein adsorption is influenced by many surface properties such as surface wettability, surface chemical composition and surface nanometre-scale morphology. Surfactant adsorption is a similar phenomenon, but utilising surfactant molecules in the place of proteins. Adsorption chillers Combining an adsorbent with a refrigerant, adsorption chillers use heat to provide a cooling effect. This heat, in the form of hot water, may come from any number of industrial sources including waste heat from industrial processes, prime heat from solar thermal installations or from the exhaust or water jacket heat of a piston engine or turbine. Although there are similarities between adsorption chillers and absorption refrigeration, the former is based on the interaction between gases and solids. The adsorption chamber of the chiller is filled with a solid material (for example zeolite, silica gel, alumina, active carbon or certain types of metal salts), which in its neutral state has adsorbed the refrigerant. When heated, the solid desorbs (releases) refrigerant vapour, which subsequently is cooled and liquefied. This liquid refrigerant then provides a cooling effect at the evaporator from its enthalpy of vaporization. In the final stage the refrigerant vapour is (re)adsorbed into the solid. As an adsorption chiller requires no compressor, it is relatively quiet. Portal site mediated adsorption Portal site mediated adsorption is a model for site-selective activated gas adsorption in metallic catalytic systems that contain a variety of different adsorption sites. In such systems, low-coordination "edge and corner" defect-like sites can exhibit significantly lower adsorption enthalpies than high-coordination (basal plane) sites. As a result, these sites can serve as "portals" for very rapid adsorption to the rest of the surface. The phenomenon relies on the common "spillover" effect (described below), where certain adsorbed species exhibit high mobility on some surfaces. The model explains seemingly inconsistent observations of gas adsorption thermodynamics and kinetics in catalytic systems where surfaces can exist in a range of coordination structures, and it has been successfully applied to bimetallic catalytic systems where synergistic activity is observed. In contrast to pure spillover, portal site adsorption refers to surface diffusion to adjacent adsorption sites, not to non-adsorptive support surfaces. The model appears to have been first proposed for carbon monoxide on silica-supported platinum by Brandt et al. (1993). A similar, but independent model was developed by King and co-workers to describe hydrogen adsorption on silica-supported alkali promoted ruthenium, silver-ruthenium and copper-ruthenium bimetallic catalysts. The same group applied the model to CO hydrogenation (Fischer–Tropsch synthesis). Zupanc et al. (2002) subsequently confirmed the same model for hydrogen adsorption on magnesia-supported caesium-ruthenium bimetallic catalysts. Trens et al. (2009) have similarly described CO surface diffusion on carbon-supported Pt particles of varying morphology. Adsorption spillover In the case catalytic or adsorbent systems where a metal species is dispersed upon a support (or carrier) material (often quasi-inert oxides, such as alumina or silica), it is possible for an adsorptive species to indirectly adsorb to the support surface under conditions where such adsorption is thermodynamically unfavorable. The presence of the metal serves as a lower-energy pathway for gaseous species to first adsorb to the metal and then diffuse on the support surface. This is possible because the adsorbed species attains a lower energy state once it has adsorbed to the metal, thus lowering the activation barrier between the gas phase species and the support-adsorbed species. Hydrogen spillover is the most common example of an adsorptive spillover. In the case of hydrogen, adsorption is most often accompanied with dissociation of molecular hydrogen (H2) to atomic hydrogen (H), followed by spillover of the hydrogen atoms present. The spillover effect has been used to explain many observations in heterogeneous catalysis and adsorption. Polymer adsorption Adsorption of molecules onto polymer surfaces is central to a number of applications, including development of non-stick coatings and in various biomedical devices. Polymers may also be adsorbed to surfaces through polyelectrolyte adsorption. In viruses Adsorption is the first step in the viral life cycle. The next steps are penetration, uncoating, synthesis (transcription if needed, and translation), and release. The virus replication cycle, in this respect, is similar for all types of viruses. Factors such as transcription may or may not be needed if the virus is able to integrate its genomic information in the cell's nucleus, or if the virus can replicate itself directly within the cell's cytoplasm. In popular culture The game of Tetris is a puzzle game in which blocks of 4 are adsorbed onto a surface during game play. Scientists have used Tetris blocks "as a proxy for molecules with a complex shape" and their "adsorption on a flat surface" for studying the thermodynamics of nanoparticles.	Physical sciences	Other separations	Chemistry
207619	https://en.wikipedia.org/wiki/Elliptical%20galaxy	Elliptical galaxy	An elliptical galaxy is a type of galaxy with an approximately ellipsoidal shape and a smooth, nearly featureless image. They are one of the three main classes of galaxy described by Edwin Hubble in his Hubble sequence and 1936 work The Realm of the Nebulae, along with spiral and lenticular galaxies. Elliptical (E) galaxies are, together with lenticular galaxies (S0) with their large-scale disks, and ES galaxies with their intermediate scale disks, a subset of the "early-type" galaxy population. Most elliptical galaxies are composed of older, low-mass stars, with a sparse interstellar medium, and they tend to be surrounded by large numbers of globular clusters. Star formation activity in elliptical galaxies is typically minimal; they may, however, undergo brief periods of star formation when merging with other galaxies. Elliptical galaxies are believed to make up approximately 10–15% of galaxies in the Virgo Supercluster, and they are not the dominant type of galaxy in the universe overall. They are preferentially found close to the centers of galaxy clusters. Elliptical galaxies range in size from dwarf ellipticals with tens of millions of stars, to supergiants of over one hundred trillion stars that dominate their galaxy clusters. Originally, Edwin Hubble hypothesized that elliptical galaxies evolved into spiral galaxies, which was later discovered to be false, although the accretion of gas and smaller galaxies may build a disk around a pre-existing ellipsoidal structure. Stars found inside of elliptical galaxies are on average much older than stars found in spiral galaxies. Examples 3C 244.1 M49 (NGC 4472) M59 (NGC 4621) M60 (NGC 4649) M87 (NGC 4486), whose supermassive black hole was the first black hole to be imaged by the Event Horizon Telescope. M89 (NGC 4552) M105 (NGC 3379) NGC 4697, part of the NGC 4697 Group ESO 383-76, one of the largest galaxies known. IC 1101, the central galaxy of Abell 2029 Hercules A, supergiant elliptical galaxy Maffei 1, the closest giant elliptical galaxy CGCG 049-033, known for having the longest galactic jet discovered Centaurus A (NGC 5128), an elliptical/lenticular radio galaxy with peculiar morphology and unusual dust lanes NeVe 1, the source of the Ophiuchus Supercluster eruption, the most powerful astronomical event known M86 (4406) an elliptical or lenticular galaxy in the constellation Virgo General characteristics Elliptical galaxies are characterized by several properties that make them distinct from other classes of galaxy. They are spherical or ovoid masses of stars, starved of star-making gases. Furthermore, there is very little interstellar matter (neither gas nor dust), which results in low rates of star formation, few open star clusters, and few young stars; rather elliptical galaxies are dominated by old stellar populations, giving them red colors. Large elliptical galaxies typically have an extensive system of globular clusters. They generally have two distinct populations of globular clusters: one that is redder and metal-rich, and another that is bluer and metal-poor. The dynamical properties of elliptical galaxies and the bulges of disk galaxies are similar, suggesting that they may be formed by the same physical processes, although this remains controversial. The luminosity profiles of both elliptical galaxies and bulges are well fit by Sersic's law, and a range of scaling relations between the elliptical galaxies' structural parameters unify the population. Every massive elliptical galaxy contains a supermassive black hole at its center. Observations of 46 elliptical galaxies, 20 classical bulges, and 22 pseudobulges show that each contain a black hole at the center. The mass of the black hole is tightly correlated with the mass of the galaxy, evidenced through correlations such as the M–sigma relation which relates the velocity dispersion of the surrounding stars to the mass of the black hole at the center. Elliptical galaxies are preferentially found in galaxy clusters and in compact groups of galaxies. Unlike flat spiral galaxies with organization and structure, elliptical galaxies are more three-dimensional, without much structure, and their stars are in somewhat random orbits around the center. Sizes and shapes The largest galaxies are supergiant ellipticals, or type-cD galaxies. Elliptical galaxies vary greatly in both size and mass with diameters ranging from 3,000 light years to more than 700,000 light years, and masses from 105 to nearly 1013 solar masses. This range is much broader for this galaxy type than for any other. The smallest, the dwarf elliptical galaxies, may be no larger than a typical globular cluster, but contain a considerable amount of dark matter not present in clusters. Most of these small galaxies may not be related to other ellipticals. The Hubble classification of elliptical galaxies contains an integer that describes how elongated the galaxy image is. The classification is determined by the ratio of the major (a) to the minor (b) axes of the galaxy's isophotes: Thus for a spherical galaxy with a equal to b, the number is 0, and the Hubble type is E0. While the limit in the literature is about E7, it has been known since 1966 that the E4 to E7 galaxies are misclassified lenticular galaxies with disks inclined at different angles to our line of sight. This has been confirmed through spectral observations revealing the rotation of their stellar disks. Hubble recognized that his shape classification depends both on the intrinsic shape of the galaxy, as well as the angle with which the galaxy is observed. Hence, some galaxies with Hubble type E0 are actually elongated. It is sometimes said that there are two physical types of ellipticals: the giant ellipticals with slightly "boxy"-shaped isophotes, whose shapes result from random motion which is greater in some directions than in others (anisotropic random motion); and the "disky" normal and dwarf ellipticals, which contain disks. This is, however, an abuse of the nomenclature, as there are two types of early-type galaxy, those with disks and those without. Given the existence of ES galaxies with intermediate-scale disks, it is reasonable to expect that there is a continuity from E to ES, and onto the S0 galaxies with their large-scale stellar disks that dominate the light at large radii. Dwarf spheroidal galaxies appear to be a distinct class: their properties are more similar to those of irregulars and late spiral-type galaxies. At the large end of the elliptical spectrum, there is further division, beyond Hubble's classification. Beyond gE giant ellipticals, lies D-galaxies and cD-galaxies. These are similar to their smaller brethren, but more diffuse, with large haloes that may as much belong to the galaxy cluster within which they reside than the centrally-located giant galaxy. Star formation In recent years, evidence has shown that a reasonable proportion (~25%) of early-type (E, ES and S0) galaxies have residual gas reservoirs and low-level star formation. Herschel Space Observatory researchers have speculated that the central black holes in elliptical galaxies keep the gas from cooling enough for star formation.	Physical sciences	Galaxy morphological classification	null
207620	https://en.wikipedia.org/wiki/Spiral%20galaxy	Spiral galaxy	Spiral galaxies form a class of galaxy originally described by Edwin Hubble in his 1936 work The Realm of the Nebulae and, as such, form part of the Hubble sequence. Most spiral galaxies consist of a flat, rotating disk containing stars, gas and dust, and a central concentration of stars known as the bulge. These are often surrounded by a much fainter halo of stars, many of which reside in globular clusters. Spiral galaxies are named by their spiral structures that extend from the center into the galactic disc. The spiral arms are sites of ongoing star formation and are brighter than the surrounding disc because of the young, hot OB stars that inhabit them. Roughly two-thirds of all spirals are observed to have an additional component in the form of a bar-like structure, extending from the central bulge, at the ends of which the spiral arms begin. The proportion of barred spirals relative to barless spirals has likely changed over the history of the universe, with only about 10% containing bars about 8 billion years ago, to roughly a quarter 2.5 billion years ago, until present, where over two-thirds of the galaxies in the visible universe (Hubble volume) have bars. The Milky Way is a barred spiral, although the bar itself is difficult to observe from Earth's current position within the galactic disc. The most convincing evidence for the stars forming a bar in the Galactic Center comes from several recent surveys, including the Spitzer Space Telescope. Together with irregular galaxies, spiral galaxies make up approximately 60% of galaxies in today's universe. They are mostly found in low-density regions and are rare in the centers of galaxy clusters. Structure Spiral galaxies may consist of several distinct components: A flat, rotating disc of stars and interstellar matter of which spiral arms are prominent components A central stellar bulge of mainly older stars, which resembles an elliptical galaxy A bar-shaped distribution of stars A near-spherical halo of stars, including many in globular clusters A supermassive black hole at the very center of the bulge A near-spherical dark matter halo The relative importance, in terms of mass, brightness and size, of the different components varies from galaxy to galaxy. Spiral arms Spiral arms are regions of stars that extend from the center of barred and unbarred spiral galaxies. These long, thin regions resemble a spiral and thus give spiral galaxies their name. Naturally, different classifications of spiral galaxies have distinct arm-structures. Sc and SBc galaxies, for instance, have very "loose" arms, whereas Sa and SBa galaxies have tightly wrapped arms (with reference to the Hubble sequence). Either way, spiral arms contain many young, blue stars (due to the high mass density and the high rate of star formation), which make the arms so bright. Bulge A bulge is a large, tightly packed group of stars. The term refers to the central group of stars found in most spiral galaxies, often defined as the excess of stellar light above the inward extrapolation of the outer (exponential) disk light. Using the Hubble classification, the bulge of Sa galaxies is usually composed of Population II stars, which are old, red stars with low metal content. Further, the bulge of Sa and SBa galaxies tends to be large. In contrast, the bulges of Sc and SBc galaxies are much smaller and are composed of young, blue Population I stars. Some bulges have similar properties to those of elliptical galaxies (scaled down to lower mass and luminosity); others simply appear as higher density centers of disks, with properties similar to disk galaxies. Many bulges are thought to host a supermassive black hole at their centers. In our own galaxy, for instance, the object called Sagittarius A* is a supermassive black hole. There are many lines of evidence for the existence of black holes in spiral galaxy centers, including the presence of active nuclei in some spiral galaxies, and dynamical measurements that find large compact central masses in galaxies such as Messier 106. Bar Bar-shaped elongations of stars are observed in roughly two-thirds of all spiral galaxies. Their presence may be either strong or weak. In edge-on spiral (and lenticular) galaxies, the presence of the bar can sometimes be discerned by the out-of-plane X-shaped or (peanut shell)-shaped structures which typically have a maximum visibility at half the length of the in-plane bar. Spheroid The bulk of the stars in a spiral galaxy are located either close to a single plane (the galactic plane) in more or less conventional circular orbits around the center of the galaxy (the Galactic Center), or in a spheroidal galactic bulge around the galactic core. However, some stars inhabit a spheroidal halo or galactic spheroid, a type of galactic halo. The orbital behaviour of these stars is disputed, but they may exhibit retrograde and/or highly inclined orbits, or not move in regular orbits at all. Halo stars may be acquired from small galaxies which fall into and merge with the spiral galaxy—for example, the Sagittarius Dwarf Spheroidal Galaxy is in the process of merging with the Milky Way and observations show that some stars in the halo of the Milky Way have been acquired from it. Unlike the galactic disc, the halo seems to be free of dust, and in further contrast, stars in the galactic halo are of Population II, much older and with much lower metallicity than their Population I cousins in the galactic disc (but similar to those in the galactic bulge). The galactic halo also contains many globular clusters. The motion of halo stars does bring them through the disc on occasion, and a number of small red dwarfs close to the Sun are thought to belong to the galactic halo, for example Kapteyn's Star and Groombridge 1830. Due to their irregular movement around the center of the galaxy, these stars often display unusually high proper motion. Oldest spiral galaxies BRI 1335-0417 is the oldest and most distant known spiral galaxy, as of 2024. The galaxy has a redshift of 4.4, meaning its light took 12.4 billion years to reach Earth. The oldest grand design spiral galaxy on file is BX442. At eleven billion years old, it is more than two billion years older than any previous discovery. Researchers believe the galaxy's shape is caused by the gravitational influence of a companion dwarf galaxy. Computer models based on that assumption indicate that BX442's spiral structure will last about 100 million years. The oldest multi-arm spiral galaxy, as of 2022, is A2744-DSG-z3. Its redshift is z=3.059, which corresponds to 11.5 billion light years to Earth. A1689B11 is an extremely old spiral galaxy located in the Abell 1689 galaxy cluster in the Virgo constellation. A1689B11 is 11 billion light years from the Earth, forming 2.6 billion years after the Big Bang. Related In June 2019, citizen scientists through Galaxy Zoo reported that the usual Hubble classification, particularly concerning spiral galaxies, may not be supported, and may need updating. Origin of the spiral structure History The pioneer of studies of the rotation of the Galaxy and the formation of the spiral arms was Bertil Lindblad in 1925. He realized that the idea of stars arranged permanently in a spiral shape was untenable. Since the angular speed of rotation of the galactic disk varies with distance from the centre of the galaxy (via a standard solar system type of gravitational model), a radial arm (like a spoke) would quickly become curved as the galaxy rotates. The arm would, after a few galactic rotations, become increasingly curved and wind around the galaxy ever tighter. This is called the winding problem. Measurements in the late 1960s showed that the orbital velocity of stars in spiral galaxies with respect to their distance from the galactic center is indeed higher than expected from Newtonian dynamics but still cannot explain the stability of the spiral structure. Since the 1970s, there have been two leading hypotheses or models for the spiral structures of galaxies: star formation caused by density waves in the galactic disk of the galaxy. the stochastic self-propagating star formation model (SSPSF model) – star formation caused by shock waves in the interstellar medium. The shock waves are caused by the stellar winds and supernovae from recent previous star formation, leading to self-propagating and self-sustaining star formation. Spiral structure then arises from differential rotation of the galaxy's disk. These different hypotheses are not mutually exclusive, as they may explain different types of spiral arms. Density wave model Bertil Lindblad proposed that the arms represent regions of enhanced density (density waves) that rotate more slowly than the galaxy's stars and gas. As gas enters a density wave, it gets squeezed and makes new stars, some of which are short-lived blue stars that light the arms. Historical theory of Lin and Shu The first acceptable theory for the spiral structure was devised by C. C. Lin and Frank Shu in 1964, attempting to explain the large-scale structure of spirals in terms of a small-amplitude wave propagating with fixed angular velocity, that revolves around the galaxy at a speed different from that of the galaxy's gas and stars. They suggested that the spiral arms were manifestations of spiral density waves – they assumed that the stars travel in slightly elliptical orbits, and that the orientations of their orbits is correlated i.e. the ellipses vary in their orientation (one to another) in a smooth way with increasing distance from the galactic center. This is illustrated in the diagram to the right. It is clear that the elliptical orbits come close together in certain areas to give the effect of arms. Stars therefore do not remain forever in the position that we now see them in, but pass through the arms as they travel in their orbits. Star formation caused by density waves The following hypotheses exist for star formation caused by density waves: As gas clouds move into the density wave, the local mass density increases. Since the criteria for cloud collapse (the Jeans instability) depends on density, a higher density makes it more likely for clouds to collapse and form stars. As the compression wave goes through, it triggers star formation on the leading edge of the spiral arms. As clouds get swept up by the spiral arms, they collide with one another and drive shock waves through the gas, which in turn causes the gas to collapse and form stars. More young stars in spiral arms Spiral arms appear visually brighter because they contain both young stars and more massive and luminous stars than the rest of the galaxy. As massive stars evolve far more quickly, their demise tends to leave a darker background of fainter stars immediately behind the density waves. This make the density waves much more prominent. Spiral arms simply appear to pass through the older established stars as they travel in their galactic orbits, so they also do not necessarily follow the arms. As stars move through an arm, the space velocity of each stellar system is modified by the gravitational force of the local higher density. Also the newly created stars do not remain forever fixed in the position within the spiral arms, where the average space velocity returns to normal after the stars depart on the other side of the arm. Gravitationally aligned orbits Charles Francis and Erik Anderson showed from observations of motions of over 20,000 local stars (within 300 parsecs) that stars do move along spiral arms, and described how mutual gravity between stars causes orbits to align on logarithmic spirals. When the theory is applied to gas, collisions between gas clouds generate the molecular clouds in which new stars form, and evolution towards grand-design bisymmetric spirals is explained. Distribution of stars in spirals The stars in spirals are distributed in thin disks radial with intensity profiles such that with being the disk scale-length; is the central value; it is useful to define: as the size of the stellar disk, whose luminosity is . The spiral galaxies light profiles, in terms of the coordinate , do not depend on galaxy luminosity. Spiral nebula Before it was understood that spiral galaxies existed outside of our Milky Way galaxy, they were often referred to as spiral nebulae, due to Lord Rosse, whose telescope Leviathan was the first to reveal the spiral structure of galaxies. In 1845 he discovered the spiral structure of M51, a galaxy nicknamed later as the "Whirlpool Galaxy", and his drawings of it closely resemble modern photographs. In 1846 and in 1849 Lord Rosse identified similar pattern in Messier 99 and Messier 33 respectively. In 1850 he made the first drawing of Andromeda Galaxy's spiral structure. In 1852 Stephen Alexander supposed that Milky Way is also a spiral nebula. The question of whether such objects were separate galaxies independent of the Milky Way, or a type of nebula existing within our own galaxy, was the subject of the Great Debate of 1920, between Heber Curtis of Lick Observatory and Harlow Shapley of Mount Wilson Observatory. Beginning in 1923, Edwin Hubble observed Cepheid variables in several spiral nebulae, including the so-called "Andromeda Nebula", proving that they are, in fact, entire galaxies outside our own. The term spiral nebula has since fallen out of use. Milky Way The Milky Way was once considered an ordinary spiral galaxy. Astronomers first began to suspect that the Milky Way is a barred spiral galaxy in the 1960s. Their suspicions were confirmed by Spitzer Space Telescope observations in 2005, which showed that the Milky Way's central bar is larger than what was previously suspected. Famous examples	Physical sciences	Galaxy morphological classification	null
207621	https://en.wikipedia.org/wiki/Irregular%20galaxy	Irregular galaxy	An irregular galaxy is a galaxy that does not have a distinct regular shape, unlike a spiral or an elliptical galaxy. Irregular galaxies do not fall into any of the regular classes of the Hubble sequence, and they are often chaotic in appearance, with neither a nuclear bulge nor any trace of spiral arm structure. Collectively they are thought to make up about a quarter of all galaxies. Some irregular galaxies were once spiral or elliptical galaxies but were deformed by an uneven external gravitational force. Irregular galaxies may contain abundant amounts of gas and dust. This is not necessarily true for dwarf irregulars. Irregular galaxies are commonly small, about one tenth the mass of the Milky Way galaxy, though there are also unusual cases of large irregulars like UGC 6697. Due to their small sizes, they are prone to environmental effects like colliding with large galaxies and intergalactic clouds. Types There are three major types of irregular galaxies: An Irr-I galaxy (Irr I) is an irregular galaxy that features some structure but not enough to place it cleanly into the Hubble sequence. Subtypes with some spiral structure are called Sm galaxies Subtypes without spiral structure are called Im galaxies. An Irr-II galaxy (Irr II) is an irregular galaxy that does not appear to feature any structure that can place it into the Hubble sequence. A dI-galaxy (or dIrr) is a dwarf irregular galaxy. This type of galaxy is now thought to be important to understand the overall evolution of galaxies, as they tend to have a low level of metallicity and relatively high levels of gas, and are thought to be similar to the earliest galaxies that populated the Universe. They may represent a local (and therefore more recent) version of the faint blue galaxies known to exist in deep field galaxy surveys. Some of the irregular galaxies, especially of the Magellanic type, are small spiral galaxies that are being distorted by the gravity of a larger neighbor. Magellanic Clouds The Magellanic Cloud galaxies were once classified as irregular galaxies. The Large Magellanic Cloud has since been re-classified as type SBm (barred Magellanic spiral). The Small Magellanic Cloud remains classified as an irregular galaxy of type Im under current galaxy morphological classification, although it does contain a bar structure. Gallery	Physical sciences	Galaxy morphological classification	null
207685	https://en.wikipedia.org/wiki/Cassiopeia%20%28constellation%29	Cassiopeia (constellation)	Cassiopeia () is a constellation and asterism in the northern sky named after the vain queen Cassiopeia, mother of Andromeda, in Greek mythology, who boasted about her unrivaled beauty. Cassiopeia was one of the 48 constellations listed by the 2nd-century Greek astronomer Ptolemy, and it remains one of the 88 modern constellations today. It is easily recognizable due to its distinctive 'W' shape, formed by five bright stars. Cassiopeia is located in the northern sky and from latitudes above 34°N it is visible year-round. In the (sub)tropics it can be seen at its clearest from September to early November, and at low southern, tropical, latitudes of less than 25°S it can be seen, seasonally, low in the North. At magnitude 2.2, Alpha Cassiopeiae, or Schedar, is the brightest star in Cassiopeia. The constellation hosts some of the most luminous stars known, including the yellow hypergiants Rho Cassiopeiae and V509 Cassiopeiae and white hypergiant 6 Cassiopeiae. In 1572, Tycho Brahe's supernova flared brightly in Cassiopeia. Cassiopeia A is a supernova remnant and the brightest extrasolar radio source in the sky at frequencies above 1 GHz. Fourteen star systems have been found to have exoplanets, one of whichHD 219134is thought to host six planets. A rich section of the Milky Way runs through Cassiopeia, containing a number of open clusters, young luminous galactic disc stars, and nebulae. IC 10 is an irregular galaxy that is the closest known starburst galaxy and the only one in the Local Group of galaxies. Mythology The constellation is named after Cassiopeia, the queen of Aethiopia. Cassiopeia was the wife of King Cepheus of Aethiopia and mother of Princess Andromeda. Cepheus and Cassiopeia were placed next to each other among the stars, along with Andromeda. She was placed in the sky as a punishment after enraging Poseidon with the boast that her daughter Andromeda was more beautiful than the Nereids or, alternatively, that she herself was more beautiful than the sea nymphs. She was forced to wheel around the north celestial pole on her throne, spending half of her time clinging to it so she does not fall off, and Poseidon decreed that Andromeda should be bound to a rock as prey for the monster Cetus. Andromeda was then rescued by the hero Perseus, whom she later married. Cassiopeia has been variously portrayed throughout her history as a constellation. In Persia, she was drawn by al-Sufi as a queen holding a staff with a crescent moon in her right hand, wearing a crown, as well as a two-humped camel. In France, she was portrayed as having a marble throne and a palm leaf in her left hand, holding her robe in her right hand. This depiction is from Augustin Royer's 1679 atlas. In Chinese astronomy, the stars forming the constellation Cassiopeia are found among three areas: the Purple Forbidden enclosure (紫微垣, Zǐ Wēi Yuán), the Black Tortoise of the North (北方玄武, Běi Fāng Xuán Wǔ), and the White Tiger of the West (西方白虎, Xī Fāng Bái Hǔ). The Chinese astronomers saw several figures in what is modern-day Cassiopeia. Kappa, Eta, and Mu Cassiopeiae formed a constellation called the Bridge of the Kings; when seen along with Alpha and Beta Cassiopeiae, they formed the great chariot Wang-Liang. The charioteer's whip was represented by Gamma Cassiopeiae, sometimes called "Tsih", the Chinese word for "whip". In Hindu Mythology, Cassiopeia was associated with the mythological figure Sharmishtha – the daughter of the great Devil (Daitya) King Vrishparva and a friend to Devayani (Andromeda). In Welsh Mythology Llys Dôn (literally "The Court of Dôn") is the traditional Welsh name for the constellation. At least three of Dôn's children also have astronomical associations: Caer Gwydion ("The fortress of Gwydion") is the traditional Welsh name for the Milky Way, and Caer Arianrhod ("The Fortress of Arianrhod") being the constellation of Corona Borealis. In the 17th century, various Biblical figures were depicted in the stars of Cassiopeia. These included Bathsheba, Solomon's mother; Deborah, an Old Testament prophet; and Mary Magdalene, a follower of Jesus. A figure called the "Tinted Hand" also appeared in the stars of Cassiopeia in some Arab atlases. This is variously said to represent a woman's hand dyed red with henna, as well as the bloodied hand of Muhammad's daughter Fatima. The hand is made up of the stars α Cas, β Cas, γ Cas, δ Cas, ε Cas, and η Cas. The arm is made up of the stars α Per, γ Per, δ Per, ε Per, η Per, and ν Per. Another Arab constellation that incorporated the stars of Cassiopeia was the Camel. Its head was composed of Lambda, Kappa, Iota, and Phi Andromedae; its hump was Beta Cassiopeiae; its body was the rest of Cassiopeia, and the legs were composed of stars in Perseus and Andromeda. Other cultures see a hand or moose antlers in the pattern. These include the Sámi, for whom the W of Cassiopeia forms an elk antler. The Chukchi of Siberia similarly saw the five main stars as five reindeer stags. The people of the Marshall Islands saw Cassiopeia as part of a great porpoise constellation. The main stars of Cassiopeia make its tail, Andromeda and Triangulum form its body, and Aries makes its head. In Hawaii, Alpha, Beta, and Gamma Cassiopeiae were named. Alpha Cassiopeiae was called Poloahilani, Beta Cassiopeiae was called Polula, and Gamma Cassiopeiae was called Mulehu. The people of Pukapuka saw the figure of Cassiopeia as a distinct constellation called Na Taki-tolu-a-Mataliki. Characteristics Cassiopeia had a supernova, Cassiopeia A, SN 1572. Covering 598.4 square degrees and hence 1.451% of the sky, Cassiopeia ranks 25th of the 88 constellations in area. It is bordered by Cepheus to the north and west, Andromeda to the south and west, Perseus to the southeast and Camelopardalis to the east, and also shares a short border with Lacerta to the west. The three-letter abbreviation for the constellation, as adopted by the International Astronomical Union in 1922, is "Cas". The official constellation boundaries, as set by Belgian astronomer Eugène Delporte in 1930, are defined by a polygon of 30 segments. In the equatorial coordinate system, the right ascension coordinates of these borders lie between and , while the declination coordinates are between 77.69° and 46.68°. Its position in the Northern Celestial Hemisphere means that the whole constellation is visible to observers north of 12°S. High in the northern sky, it is circumpolar (that is, it never sets in the night sky) to viewers in the British Isles, Canada and the northern United States. Features Stars The German cartographer Johann Bayer used the Greek letters Alpha through Omega, and then A and B, to label the most prominent 26 stars in the constellation. Upsilon was later found to be two stars and labelled Upsilon1 and Upsilon2 by John Flamsteed. B Cassiopeiae was in fact the supernova known as Tycho's Supernova. Within the constellation's borders, there are 157 stars brighter than or equal to apparent magnitude 6.5. 'W' asterism The five brightest stars of CassiopeiaAlpha, Beta, Gamma, Delta, and Epsilon Cassiopeiaeform the characteristic W-shaped asterism. All five are prominent naked eye stars, three are noticeably variable, and a fourth is a suspected low amplitude variable. The asterism is oriented as a W when below Polaris during northern spring and summer nights. In northern winter, and when seen from southern latitudes, it is "above" Polaris (i.e. closer to the zenith) and the W appears inverted. Alpha Cassiopeiae, traditionally called Schedar (from the Arabic Al Sadr, "the breast"), is commonly mistaken as a four-star system, but is actually a single star with three physically distant optical components. The primary dominates, an orange-hued giant of magnitude 2.2, from Earth. With a luminosity of around 771 times that of the Sun, it has swollen and cooled after exhausting its core hydrogen over its 100 to 200 million-year lifespan, spending much of it as a blue-white B-type main-sequence star. Magnitude 8.9 yellow dwarf companion (B) is widely separated; companions (C and D) are closer and magnitudes 13 and 14 respectively. Beta Cassiopeiae, or Caph (meaning "hand"), is a white-hued star of magnitude 2.3, from Earth. Around 1.2 billion years old, it has used up its core hydrogen and begun expanding and cooling off the main sequence. It is around 1.9 times as massive as the Sun, and around 21.3 times as luminous. Rotating at about 92% of its critical speed, Caph completes a full rotation every 1.12 days. This is giving the star an oblate spheroid shape with an equatorial bulge that is 24% larger than the polar radius. It is a Delta Scuti variable with a small amplitude and period of 2.5 hours. Gamma Cassiopeiae is the prototype Gamma Cassiopeiae variable star, a type of variable star that has a variable disc of material flung off by the high rotation rate of the star. Gamma Cassiopeiae has a minimum magnitude of 3.0 and a maximum magnitude of 1.6, but is generally near magnitude 2.2, with unpredictable fades and brightenings. It is a spectroscopic binary, with an orbital period of 203.59 days and a companion with a calculated mass about the same as the Sun. However, no direct evidence of this companion has been found, leading to speculation that it might be a white dwarf or other degenerate star. It is from Earth. Delta Cassiopeiae, also known as Ruchbah or Rukbat, meaning "knee," is a possible Algol-type eclipsing binary star with a maximum brightness of magnitude 2.7. It has been reported to show eclipses of less than 0.1 magnitudes with a period of 2 years and 1 month., but this has never been confirmed. It is from Earth. Delta Cassiopeiae was the star used by Jean Picard in 1669 to establish the length on the Earth's surface corresponding to one degree, and hence the radius of the Earth. Epsilon Cassiopeiae has an apparent magnitude of 3.3. Located from Earth, it is a hot blue-white star of spectral type B3 III with a surface temperature of . It is 6.5 times as massive and 4.2 times as wide as the Sun, and belongs to a class of stars known as Be stars—rapidly spinning stars that throw off a ring or shell of matter. Fainter stars The next seven brightest stars in Cassiopeia are also all confirmed or suspected variable stars, including 50 Cassiopeiae which was not given a Greek letter by Bayer and is a suspected variable with a very small amplitude. Zeta Cassiopeiae (Fulu) is a suspected slowly pulsating B-type star. Kappa Cassiopeiae is a blue supergiant of spectral type BC0.7Ia that is some 302,000 times as luminous as the Sun and has 33 times its diameter. It is a runaway star, moving at around 2.5 million mph relative to its neighbors (1,100 kilometers per second). Its magnetic field and wind of particles creates a visible bow shock 4 light-years ahead of it, colliding with the diffuse, and usually invisible, interstellar gas and dust. The dimensions of the bow shock are vast: around 12 light-years long and 1.8 light-years wide. Theta Cassiopeiae, named Marfak, is a suspected variable star whose brightness changes by less than a tenth of a magnitude. Iota Cassiopeiae is a triple star 142 light-years from Earth. The primary is a white-hued star of magnitude 4.5 and an α2 Canum Venaticorum variable, the secondary is a yellow-hued star of magnitude 6.9, and the tertiary is a star of magnitude 8.4. The primary and secondary are close together but the primary and tertiary are widely separated. Omicron Cassiopeiae is a triple star and the primary is another γ Cassiopeiae variable. Eta Cassiopeiae (Achird) is the nearest star in Cassiopeia, located some 19.13 ly away. It is a binary star made up of a G-dwarf star similar to the Sun, of apparent magnitude 3.45, and a K-dwarf star of apparent magnitude 7.51. It is visually located between Alpha and Gamma Cassiopeiae and are one of the nearest G-type / K-type stars to Earth. Sigma Cassiopeiae is a binary star 1500 light-years from Earth. It has a green-hued primary of magnitude 5.0 and a blue-hued secondary of magnitude 7.3. Psi Cassiopeiae is a triple star 193 light-years from Earth. The primary is an orange-hued giant star of magnitude 4.7 and the secondary is a close pair of stars that appears to be of magnitude 9.0. Rho Cassiopeiae is a semi-regular pulsating variable yellow hypergiant, and is among the most luminous stars in the galaxy at approximately . It has a minimum magnitude of 6.2 and a maximum magnitude of 4.1; its period is approximately 320 days. It has around 450 times the Sun's diameter and 17 times its mass, having begun life 45 times as massive as the Sun. Rho Cassiopeiae is about 10,000 light-years from Earth. Cassiopeia includes V509 Cassiopeiae, a second example of the extremely rare yellow hypergiants, which is around 400,000 times as luminous as the Sun and 14 times as massive, as well as 6 Cassiopeiae which is a hotter white hypergiant. It also hosts the red supergiant PZ Cassiopeiae, which is one of the largest known stars with an estimate of and is also a semiregular variable. Between 240,000 and 270,000 times as luminous as the Sun, it is around 9,160 light-years distant from Earth. AO Cassiopeiae is a binary system composed of an O8 main sequence star and an O9.2 bright giant that respectively weigh anywhere between 20.30 and 57.75 times and 14.8 and 31.73 times the mass of the Sun. The two massive stars are so close to each other they distort each other into egg-shapes. Tycho Brahe's supernova was visible within Cassiopeia, and the star Tycho G is suspected of being the donor of the material that triggered that explosion. Deep-sky objects A rich section of the Milky Way runs through Cassiopeia, stretching from Perseus towards Cygnus, and it contains a number of open clusters, young luminous galactic disc stars, and nebulae. The Heart Nebula and the Soul Nebula are two neighboring emission nebulae about 7,500 light-years away. Two Messier objects, M52 (NGC 7654) and M103 (NGC 581), are located in Cassiopeia; both are open clusters. M52, once described as a "kidney-shaped" cluster, contains approximately 100 stars and is from Earth. Its most prominent member is an orange-hued star of magnitude 8.0 near the cluster's edge. M103 is far poorer than M52, with only about 25 stars included. It is also more distant, between 8000 and 9500 light-years from Earth. Its most prominent member is actually a closer, superimposed double star; it consists of a 7th-magnitude primary and 10th-magnitude secondary. The other prominent open clusters in Cassiopeia are NGC 457 and NGC 663, both of which have about 80 stars. NGC 457 is looser, and its brightest member is Phi Cassiopeiae, a white-hued supergiant star of magnitude 5.0. However, it is uncertain whether Phi Cassiopeiae is part of the open cluster or not. The stars of NGC 457, arrayed in chains, are approximately 10,000 light-years from Earth. NGC 663 is both closer, at 8200 light-years from Earth, and larger, at 0.25 degrees in diameter. There are two supernova remnants in Cassiopeia. The first, designated 3C 10 or just Tycho's Supernova Remnant, is the aftermath of the supernova called Tycho's Star. It was observed in 1572 by Tycho Brahe and now exists as a bright object in the radio spectrum. Within the 'W' asterism formed by Cassiopeia's five major stars lies Cassiopeia A (Cas A). It is the remnant of a supernova that took place approximately 300 years ago (as observed now from Earth; it is 10,000 light-years away), and has the distinction of being the strongest radio source observable outside the Solar System. It was perhaps observed as a faint star in 1680 by John Flamsteed. It was also the subject of the first image returned by the Chandra X-Ray Observatory in the late 1990s. The shell of matter expelled from the star is moving at per second; it has a temperature of on average. NGC 457 is another open cluster in Cassiopeia, also called the E.T. Cluster, the Owl Cluster, and Caldwell 13. The cluster was discovered in 1787 by William Herschel. It has an overall magnitude of 6.4 and is approximately 10,000 light-years from Earth, lying in the Perseus Arm of the Milky Way. However, its most prominent member, the double star Phi Cassiopeiae, is far closerbetween 1000 and 4000 light-years away. NGC 457 is fairly rich; it is a Shapley class e and Trumpler class I 3 r cluster. It is concentrated towards its center and detached from the star field. It contains more than 100 stars, which vary widely in brightness. Two members of the Local Group of galaxies are in Cassiopeia. NGC 185 is a magnitude 9.2 elliptical galaxy of type E0, 2 million light-years away. Slightly dimmer and more distant NGC 147 is a magnitude 9.3 elliptical galaxy, like NGC 185 it is an elliptical of type E0; it is 2.3 million light-years from Earth. Though they do not appear in Andromeda, both dwarf galaxies are gravitationally bound to the far larger Andromeda Galaxy. IC 10 is an irregular galaxy that is the closest known starburst galaxy and the only one in the Local Group of galaxies. Cassiopeia also contains part of the closest galaxy group to our Local Group, the IC 342/Maffei Group. The galaxies Maffei 1 and Maffei 2 are located just to the south of the Heart and Soul nebulae. As a result of this location in the Zone of Avoidance, both are surprisingly faint despite both being within 10 million light-years away (Maffei 2 is below the range of most amateur telescopes). Meteor shower The December Phi Cassiopeiids are a recently discovered early December meteor shower that radiates from Cassiopeia. Phi Cassiopeiids are very slow, with an entry velocity of approximately 16.7 kilometers per second. The shower's parent body is a Jupiter family comet, though its specific identity is unknown. Namesakes USS Cassiopeia (AK-75) was a United States Navy Crater-class cargo ship named after the constellation. In Pokémon Scarlet and Violet, the villainous team, Team Star, is divided into five squads named after the brightest stars in the constellation: Segin Squad, Schedar Squad, Ruchbah Squad, Navi Squad, and Caph Squad. The group's leader uses the alias Cassiopeia. In Ni no Kuni: Wrath of the White Witch, the penultimate main antagonist and "White Witch" in question is named Queen Cassiopeia. Cassiopeia is also the name of a song by London-based band Bears in Trees. Although the lyrics of the song mainly refer to the ancient Greek woman, the album cover shows the constellation. Cassiopeia is the name of a champion in League of Legends. Her beauty and vanity mirror the character in Greek mythology. Casiopea is the name of a Japanese Jazz-Fusion group formed in 1976. The name was chosen by the guitarist's mother, solely so the band members could remember the name.	Physical sciences	Constellations	null
207754	https://en.wikipedia.org/wiki/Somatic%20cell	Somatic cell	In cellular biology, a somatic cell (), or vegetal cell, is any biological cell forming the body of a multicellular organism other than a gamete, germ cell, gametocyte or undifferentiated stem cell. Somatic cells compose the body of an organism and divide through mitosis. In contrast, gametes derive from meiosis within the germ cells of the germline and they fuse during sexual reproduction. Stem cells also can divide through mitosis, but are different from somatic in that they differentiate into diverse specialized cell types. In mammals, somatic cells make up all the internal organs, skin, bones, blood and connective tissue, while mammalian germ cells give rise to spermatozoa and ova which fuse during fertilization to produce a cell called a zygote, which divides and differentiates into the cells of an embryo. There are approximately 220 types of somatic cell in the human body. Theoretically, these cells are not germ cells (the source of gametes); they transmit their mutations, to their cellular descendants (if they have any), but not to the organism's descendants. However, in sponges, non-differentiated somatic cells form the germ line and, in Cnidaria, differentiated somatic cells are the source of the germline. Mitotic cell division is only seen in diploid somatic cells. Only some cells like germ cells take part in reproduction. Evolution As multicellularity was theorized to be evolved many times, so did sterile somatic cells. The evolution of an immortal germline producing specialized somatic cells involved the emergence of mortality, and can be viewed in its simplest version in volvocine algae. Those species with a separation between sterile somatic cells and a germline are called Weismannists. Weismannist development is relatively rare (e.g., vertebrates, arthropods, Volvox), as many species have the capacity for somatic embryogenesis (e.g., land plants, most algae, and numerous invertebrates). Genetics and chromosomes Like all cells, somatic cells contain DNA arranged in chromosomes. If a somatic cell contains chromosomes arranged in pairs, it is called diploid and the organism is called a diploid organism. The gametes of diploid organisms contain only single unpaired chromosomes and are called haploid. Each pair of chromosomes comprises one chromosome inherited from the father and one inherited from the mother. In humans, somatic cells contain 46 chromosomes organized into 23 pairs. By contrast, gametes of diploid organisms contain only half as many chromosomes. In humans, this is 23 unpaired chromosomes. When two gametes (i.e. a spermatozoon and an ovum) meet during conception, they fuse together, creating a zygote. Due to the fusion of the two gametes, a human zygote contains 46 chromosomes (i.e. 23 pairs). A large number of species have the chromosomes in their somatic cells arranged in fours ("tetraploid") or even sixes ("hexaploid"). Thus, they can have diploid or even triploid germline cells. An example of this is the modern cultivated species of wheat, Triticum aestivum L., a hexaploid species whose somatic cells contain six copies of every chromatid. The frequency of spontaneous mutations is significantly lower in advanced male germ cells than in somatic cell types from the same individual. Female germ cells also show a mutation frequency that is lower than that in corresponding somatic cells and similar to that in male germ cells. These findings appear to reflect employment of more effective mechanisms to limit the initial occurrence of spontaneous mutations in germ cells than in somatic cells. Such mechanisms likely include elevated levels of DNA repair enzymes that ameliorate most potentially mutagenic DNA damages. Cloning In recent years, the technique of cloning whole organisms has been developed in mammals, allowing almost identical genetic clones of an animal to be produced. One method of doing this is called "somatic cell nuclear transfer" and involves removing the nucleus from a somatic cell, usually a skin cell. This nucleus contains all of the genetic information needed to produce the organism it was removed from. This nucleus is then injected into an ovum of the same species which has had its own genetic material removed. The ovum now no longer needs to be fertilized, because it contains the correct amount of genetic material (a diploid number of chromosomes). In theory, the ovum can be implanted into the uterus of a same-species animal and allowed to develop. The resulting animal will be a nearly genetically identical clone to the animal from which the nucleus was taken. The only difference is caused by any mitochondrial DNA that is retained in the ovum, which is different from the cell that donated the nucleus. In practice, this technique has so far been problematic, although there have been a few high-profile successes, such as Dolly the Sheep (July 5, 1996 - February 14, 2003) and, more recently, Snuppy (April 24, 2005 - May 2015), the first cloned dog. Biobanking Somatic cells have also been collected in the practice of biobanking. The cryoconservation of animal genetic resources is a means of conserving animal genetic material in response to decreasing ecological biodiversity. As populations of living organisms fall so does their genetic diversity. This places species long-term survivability at risk. Biobanking aims to preserve biologically viable cells through long-term storage for later use. Somatic cells have been stored with the hopes that they can be reprogrammed into induced pluripotent stem cells (iPSCs), which can then differentiate into viable reproductive cells. Genetic modifications Development of biotechnology has allowed for the genetic manipulation of somatic cells, whether for the modelling of chronic disease or for the prevention of malaise conditions. Two current means of gene editing are the use of transcription activator-like effector nucleases (TALENs) or clustered regularly interspaced short palindromic repeats (CRISPR). Genetic engineering of somatic cells has resulted in some controversies, although the International Summit on Human Gene Editing has released a statement in support of genetic modification of somatic cells, as the modifications thereof are not passed on to offspring. Cellular aging In mammals a high level of repair and maintenance of cellular DNA appears to be beneficial early in life. However, some types of cell, such as those of the brain and muscle, undergo a transition from mitotic cell division to a post-mitotic (non-dividing) condition during early development, and this transition is accompanied by a reduction in DNA repair capability. This reduction may be an evolutionary adaptation permitting the diversion of cellular resources that were earlier used for DNA repair, as well as for DNA replication and cell division, to higher priority neuronal and muscular functions. An effect of these reductions is to allow increased accumulation of DNA damage likely contributing to cellular aging.	Biology and health sciences	Cell processes	Biology
207790	https://en.wikipedia.org/wiki/G%C3%B6del%20numbering	Gödel numbering	In mathematical logic, a Gödel numbering is a function that assigns to each symbol and well-formed formula of some formal language a unique natural number, called its Gödel number. Kurt Gödel developed the concept for the proof of his incompleteness theorems. () A Gödel numbering can be interpreted as an encoding in which a number is assigned to each symbol of a mathematical notation, after which a sequence of natural numbers can then represent a sequence of symbols. These sequences of natural numbers can again be represented by single natural numbers, facilitating their manipulation in formal theories of arithmetic. Since the publishing of Gödel's paper in 1931, the term "Gödel numbering" or "Gödel code" has been used to refer to more general assignments of natural numbers to mathematical objects. Simplified overview Gödel noted that each statement within a system can be represented by a natural number (its Gödel number). The significance of this was that properties of a statement—such as its truth or falsehood—would be equivalent to determining whether its Gödel number had certain properties. The numbers involved might be very large indeed, but this is not a barrier; all that matters is that such numbers can be constructed. In simple terms, Gödel devised a method by which every formula or statement that can be formulated in the system gets a unique number, in such a way that formulas and Gödel numbers can be mechanically converted back and forth. There are many ways to do this. A simple example is the way in which English is stored as a sequence of numbers in computers using ASCII. Since ASCII codes are in the range 0 to 127, it is sufficient to pad them to 3 decimal digits and then to concatenate them: The word is represented by . The logical formula is represented by . Gödel's encoding Gödel used a system based on prime factorization. He first assigned a unique natural number to each basic symbol in the formal language of arithmetic with which he was dealing. To encode an entire formula, which is a sequence of symbols, Gödel used the following system. Given a sequence of positive integers, the Gödel encoding of the sequence is the product of the first n primes raised to their corresponding values in the sequence: According to the fundamental theorem of arithmetic, any number (and, in particular, a number obtained in this way) can be uniquely factored into prime factors, so it is possible to recover the original sequence from its Gödel number (for any given number n of symbols to be encoded). Gödel specifically used this scheme at two levels: first, to encode sequences of symbols representing formulas, and second, to encode sequences of formulas representing proofs. This allowed him to show a correspondence between statements about natural numbers and statements about the provability of theorems about natural numbers, the proof's key observation (). There are more sophisticated (and more concise) ways to construct a Gödel numbering for sequences. Example In the specific Gödel numbering used by Nagel and Newman, the Gödel number for the symbol "0" is 6 and the Gödel number for the symbol "=" is 5. Thus, in their system, the Gödel number of the formula "0 = 0" is 26 × 35 × 56 = 243,000,000. Lack of uniqueness Infinitely many different Gödel numberings are possible. For example, supposing there are K basic symbols, an alternative Gödel numbering could be constructed by invertibly mapping this set of symbols (through, say, an invertible function h) to the set of digits of a bijective base-K numeral system. A formula consisting of a string of n symbols would then be mapped to the number In other words, by placing the set of K basic symbols in some fixed order, such that the -th symbol corresponds uniquely to the -th digit of a bijective base-K numeral system, each formula may serve just as the very numeral of its own Gödel number. For example, the numbering described here has K=1000. Application to formal arithmetic Recursion One may use Gödel numbering to show how functions defined by course-of-values recursion are in fact primitive recursive functions. Expressing statements and proofs by numbers Once a Gödel numbering for a formal theory is established, each inference rule of the theory can be expressed as a function on the natural numbers. If f is the Gödel mapping and r is an inference rule, then there should be some arithmetical function gr of natural numbers such that if formula C is derived from formulas A and B through an inference rule r, i.e. then This is true for the numbering Gödel used, and for any other numbering where the encoded formula can be arithmetically recovered from its Gödel number. Thus, in a formal theory such as Peano arithmetic in which one can make statements about numbers and their arithmetical relationships to each other, one can use a Gödel numbering to indirectly make statements about the theory itself. This technique allowed Gödel to prove results about the consistency and completeness properties of formal systems. Generalizations In computability theory, the term "Gödel numbering" is used in settings more general than the one described above. It can refer to: Any assignment of the elements of a formal language to natural numbers in such a way that the numbers can be manipulated by an algorithm to simulate manipulation of elements of the formal language. More generally, an assignment of elements from a countable mathematical object, such as a countable group, to natural numbers to allow algorithmic manipulation of the mathematical object. Also, the term Gödel numbering is sometimes used when the assigned "numbers" are actually strings, which is necessary when considering models of computation such as Turing machines that manipulate strings rather than numbers. Gödel sets Gödel sets are sometimes used in set theory to encode formulas, and are similar to Gödel numbers, except that one uses sets rather than numbers to do the encoding. In simple cases when one uses a hereditarily finite set to encode formulas this is essentially equivalent to the use of Gödel numbers, but somewhat easier to define because the tree structure of formulas can be modeled by the tree structure of sets. Gödel sets can also be used to encode formulas in infinitary languages.	Mathematics	Computability theory	null
207820	https://en.wikipedia.org/wiki/Compact%20object	Compact object	In astronomy, the term compact object (or compact star) refers collectively to white dwarfs, neutron stars, and black holes. It could also include exotic stars if such hypothetical, dense bodies are confirmed to exist. All compact objects have a high mass relative to their radius, giving them a very high density, compared to ordinary atomic matter. Compact objects are often the endpoints of stellar evolution and, in this respect, are also called stellar remnants. The state and type of a stellar remnant depends primarily on the mass of the star that it formed from. The ambiguous term compact object is often used when the exact nature of the star is not known, but evidence suggests that it has a very small radius compared to ordinary stars. A compact object that is not a black hole may be called a degenerate star. In June 2020, astronomers reported narrowing down the source of Fast Radio Bursts (FRBs), which may now plausibly include "compact-object mergers and magnetars arising from normal core collapse supernovae". Formation The usual endpoint of stellar evolution is the formation of a compact star. All active stars will eventually come to a point in their evolution when the outward radiation pressure from the nuclear fusions in its interior can no longer resist the ever-present gravitational forces. When this happens, the star collapses under its own weight and undergoes the process of stellar death. For most stars, this will result in the formation of a very dense and compact stellar remnant, also known as a compact star. Compact objects have no internal energy production, but will—with the exception of black holes—usually radiate for millions of years with excess heat left from the collapse itself. According to the most recent understanding, compact stars could also form during the phase separations of the early Universe following the Big Bang. Primordial origins of known compact objects have not been determined with certainty. Lifetime Although compact objects may radiate, and thus cool off and lose energy, they do not depend on high temperatures to maintain their structure, as ordinary stars do. Barring external disturbances and proton decay, they can persist virtually forever. Black holes are however generally believed to finally evaporate from Hawking radiation after trillions of years. According to our current standard models of physical cosmology, all stars will eventually evolve into cool and dark compact stars, by the time the Universe enters the so-called degenerate era in a very distant future. A somewhat wider definition of compact objects may include smaller solid objects such as planets, asteroids, and comets, but such usage is less common. There are a remarkable variety of stars and other clumps of hot matter, but all matter in the Universe must eventually end as dispersed cold particles or some form of compact stellar or substellar object, according to thermodynamics. White dwarfs The stars called white or degenerate dwarfs are made up mainly of degenerate matter; typically carbon and oxygen nuclei in a sea of degenerate electrons. White dwarfs arise from the cores of main-sequence stars and are therefore very hot when they are formed. As they cool they will redden and dim until they eventually become dark black dwarfs. White dwarfs were observed in the 19th century, but the extremely high densities and pressures they contain were not explained until the 1920s. The equation of state for degenerate matter is "soft", meaning that adding more mass will result in a smaller object. Continuing to add mass to what begins as a white dwarf, the object shrinks and the central density becomes even greater, with higher degenerate-electron energies. After the degenerate star's mass has grown sufficiently that its radius has shrunk to only a few thousand kilometers, the mass will be approaching the Chandrasekhar limit – the theoretical upper limit of the mass of a white dwarf, about 1.4 times the mass of the Sun (). If matter were removed from the center of a white dwarf and slowly compressed, electrons would first be forced to combine with nuclei, changing their protons to neutrons by inverse beta decay. The equilibrium would shift towards heavier, neutron-richer nuclei that are not stable at everyday densities. As the density increases, these nuclei become still larger and less well-bound. At a critical density of about 4 kg/m3 – called the neutron drip line – the atomic nucleus would tend to dissolve into unbound protons and neutrons. If further compressed, eventually it would reach a point where the matter is on the order of the density of an atomic nucleus – about 2 kg/m3. At that density the matter would be chiefly free neutrons, with a light scattering of protons and electrons. Neutron stars In certain binary stars containing a white dwarf, mass is transferred from the companion star onto the white dwarf, eventually pushing it over the Chandrasekhar limit. Electrons react with protons to form neutrons and thus no longer supply the necessary pressure to resist gravity, causing the star to collapse. If the center of the star is composed mostly of carbon and oxygen then such a gravitational collapse will ignite runaway fusion of the carbon and oxygen, resulting in a Type Ia supernova that entirely blows apart the star before the collapse can become irreversible. If the center is composed mostly of magnesium or heavier elements, the collapse continues. As the density further increases, the remaining electrons react with the protons to form more neutrons. The collapse continues until (at higher density) the neutrons become degenerate. A new equilibrium is possible after the star shrinks by three orders of magnitude, to a radius between 10 and 20 km. This is a neutron star. Although the first neutron star was not observed until 1967 when the first radio pulsar was discovered, neutron stars were proposed by Baade and Zwicky in 1933, only one year after the neutron was discovered in 1932. They realized that because neutron stars are so dense, the collapse of an ordinary star to a neutron star would liberate a large amount of gravitational potential energy, providing a possible explanation for supernovae. This is the explanation for supernovae of types Ib, Ic, and II. Such supernovae occur when the iron core of a massive star exceeds the Chandrasekhar limit and collapses to a neutron star. Like electrons, neutrons are fermions. They therefore provide neutron degeneracy pressure to support a neutron star against collapse. In addition, repulsive neutron-neutron interactions provide additional pressure. Like the Chandrasekhar limit for white dwarfs, there is a limiting mass for neutron stars: the Tolman–Oppenheimer–Volkoff limit, where these forces are no longer sufficient to hold up the star. As the forces in dense hadronic matter are not well understood, this limit is not known exactly but is thought to be between 2 and . If more mass accretes onto a neutron star, eventually this mass limit will be reached. What happens next is not completely clear. Black holes As more mass is accumulated, equilibrium against gravitational collapse exceeds its breaking point. Once the star's pressure is insufficient to counterbalance gravity, a catastrophic gravitational collapse occurs within milliseconds. The escape velocity at the surface, already at least light speed, quickly reaches the velocity of light. At that point no energy or matter can escape and a black hole has formed. Because all light and matter is trapped within an event horizon, a black hole appears truly black, except for the possibility of very faint Hawking radiation. It is presumed that the collapse will continue inside the event horizon. In the classical theory of general relativity, a gravitational singularity occupying no more than a point will form. There may be a new halt of the catastrophic gravitational collapse at a size comparable to the Planck length, but at these lengths there is no known theory of gravity to predict what will happen. Adding any extra mass to the black hole will cause the radius of the event horizon to increase linearly with the mass of the central singularity. This will induce certain changes in the properties of the black hole, such as reducing the tidal stress near the event horizon, and reducing the gravitational field strength at the horizon. However, there will not be any further qualitative changes in the structure associated with any mass increase. Alternative black hole models Fuzzball Gravastar Dark-energy star Black star Magnetospheric eternally collapsing object Dark star Primordial black holes Exotic stars An exotic star is a hypothetical compact star composed of something other than electrons, protons, and neutrons balanced against gravitational collapse by degeneracy pressure or other quantum properties. These include strange stars (composed of strange matter) and the more speculative preon stars (composed of preons). Exotic stars are hypothetical, but observations released by the Chandra X-Ray Observatory on April 10, 2002, detected two candidate strange stars, designated RX J1856.5-3754 and 3C58, which had previously been thought to be neutron stars. Based on the known laws of physics, the former appeared much smaller and the latter much colder than they should, suggesting that they are composed of material denser than neutronium. However, these observations are met with skepticism by researchers who say the results were not conclusive. Quark stars and strange stars If neutrons are squeezed enough at a high temperature, they will decompose into their component quarks, forming what is known as a quark matter. In this case, the star will shrink further and become denser, but instead of a total collapse into a black hole, it is possible that the star may stabilize itself and survive in this state indefinitely, so long as no more mass is added. It has, to an extent, become a very large nucleon. A star in this hypothetical state is called a "quark star" or more specifically a "strange star". The pulsar 3C58 has been suggested as a possible quark star. Most neutron stars are thought to hold a core of quark matter but this has proven difficult to determine observationally. Preon stars A preon star is a proposed type of compact star made of preons, a group of hypothetical subatomic particles. Preon stars would be expected to have huge densities, exceeding 1023 kilogram per cubic meter – intermediate between quark stars and black holes. Preon stars could originate from supernova explosions or the Big Bang; however, current observations from particle accelerators speak against the existence of preons. Q stars Q stars are hypothetical compact, heavier neutron stars with an exotic state of matter where particle numbers are preserved with radii less than 1.5 times the corresponding Schwarzschild radius. Q stars are also called "gray holes". Electroweak stars An electroweak star is a theoretical type of exotic star, whereby the gravitational collapse of the star is prevented by radiation pressure resulting from electroweak burning, that is, the energy released by conversion of quarks to leptons through the electroweak force. This process occurs in a volume at the star's core approximately the size of an apple, containing about two Earth masses. Boson star A boson star is a hypothetical astronomical object that is formed out of particles called bosons (conventional stars are formed out of fermions). For this type of star to exist, there must be a stable type of boson with repulsive self-interaction. As of 2016 there is no significant evidence that such a star exists. However, it may become possible to detect them by the gravitational radiation emitted by a pair of co-orbiting boson stars. Compact relativistic objects and the generalized uncertainty principle Based on the generalized uncertainty principle (GUP), proposed by some approaches to quantum gravity such as string theory and doubly special relativity, the effect of GUP on the thermodynamic properties of compact stars with two different components has been studied recently. Tawfik et al. noted that the existence of quantum gravity correction tends to resist the collapse of stars if the GUP parameter is taking values between Planck scale and electroweak scale. Comparing with other approaches, it was found that the radii of compact stars should be smaller and increasing energy decreases the radii of the compact stars.	Physical sciences	Stellar astronomy	Astronomy
207833	https://en.wikipedia.org/wiki/Radial%20velocity	Radial velocity	The radial velocity or line-of-sight velocity of a target with respect to an observer is the rate of change of the vector displacement between the two points. It is formulated as the vector projection of the target-observer relative velocity onto the relative direction or line-of-sight (LOS) connecting the two points. The radial speed or range rate is the temporal rate of the distance or range between the two points. It is a signed scalar quantity, formulated as the scalar projection of the relative velocity vector onto the LOS direction. Equivalently, radial speed equals the norm of the radial velocity, modulo the sign. In astronomy, the point is usually taken to be the observer on Earth, so the radial velocity then denotes the speed with which the object moves away from the Earth (or approaches it, for a negative radial velocity). Formulation Given a differentiable vector defining the instantaneous relative position of a target with respect to an observer. Let the instantaneous relative velocity of the target with respect to the observer be The magnitude of the position vector is defined as in terms of the inner product The quantity range rate is the time derivative of the magnitude (norm) of , expressed as Substituting () into () Evaluating the derivative of the right-hand-side by the chain rule using () the expression becomes By reciprocity, . Defining the unit relative position vector (or LOS direction), the range rate is simply expressed as i.e., the projection of the relative velocity vector onto the LOS direction. Further defining the velocity direction , with the relative speed , we have: where the inner product is either +1 or -1, for parallel and antiparallel vectors, respectively. A singularity exists for coincident observer target, i.e., ; in this case, range rate is undefined. Applications in astronomy In astronomy, radial velocity is often measured to the first order of approximation by Doppler spectroscopy. The quantity obtained by this method may be called the barycentric radial-velocity measure or spectroscopic radial velocity. However, due to relativistic and cosmological effects over the great distances that light typically travels to reach the observer from an astronomical object, this measure cannot be accurately transformed to a geometric radial velocity without additional assumptions about the object and the space between it and the observer. By contrast, astrometric radial velocity is determined by astrometric observations (for example, a secular change in the annual parallax). Spectroscopic radial velocity Light from an object with a substantial relative radial velocity at emission will be subject to the Doppler effect, so the frequency of the light decreases for objects that were receding (redshift) and increases for objects that were approaching (blueshift). The radial velocity of a star or other luminous distant objects can be measured accurately by taking a high-resolution spectrum and comparing the measured wavelengths of known spectral lines to wavelengths from laboratory measurements. A positive radial velocity indicates the distance between the objects is or was increasing; a negative radial velocity indicates the distance between the source and observer is or was decreasing. William Huggins ventured in 1868 to estimate the radial velocity of Sirius with respect to the Sun, based on observed redshift of the star's light. In many binary stars, the orbital motion usually causes radial velocity variations of several kilometres per second (km/s). As the spectra of these stars vary due to the Doppler effect, they are called spectroscopic binaries. Radial velocity can be used to estimate the ratio of the masses of the stars, and some orbital elements, such as eccentricity and semimajor axis. The same method has also been used to detect planets around stars, in the way that the movement's measurement determines the planet's orbital period, while the resulting radial-velocity amplitude allows the calculation of the lower bound on a planet's mass using the binary mass function. Radial velocity methods alone may only reveal a lower bound, since a large planet orbiting at a very high angle to the line of sight will perturb its star radially as much as a much smaller planet with an orbital plane on the line of sight. It has been suggested that planets with high eccentricities calculated by this method may in fact be two-planet systems of circular or near-circular resonant orbit. Detection of exoplanets The radial velocity method to detect exoplanets is based on the detection of variations in the velocity of the central star, due to the changing direction of the gravitational pull from an (unseen) exoplanet as it orbits the star. When the star moves towards us, its spectrum is blueshifted, while it is redshifted when it moves away from us. By regularly looking at the spectrum of a star—and so, measuring its velocity—it can be determined if it moves periodically due to the influence of an exoplanet companion. Data reduction From the instrumental perspective, velocities are measured relative to the telescope's motion. So an important first step of the data reduction is to remove the contributions of the Earth's elliptic motion around the Sun at approximately ± 30 km/s, a monthly rotation of ± 13 m/s of the Earth around the center of gravity of the Earth-Moon system, the daily rotation of the telescope with the Earth crust around the Earth axis, which is up to ±460 m/s at the equator and proportional to the cosine of the telescope's geographic latitude, small contributions from the Earth polar motion at the level of mm/s, contributions of 230 km/s from the motion around the Galactic Center and associated proper motions. in the case of spectroscopic measurements corrections of the order of ±20 cm/s with respect to aberration. Sin i degeneracy is the impact caused by not being in the plane of the motion.	Physical sciences	Basics	Astronomy
207874	https://en.wikipedia.org/wiki/Swarm%20behaviour	Swarm behaviour	Swarm behaviour, or swarming, is a collective behaviour exhibited by entities, particularly animals, of similar size which aggregate together, perhaps milling about the same spot or perhaps moving en masse or migrating in some direction. It is a highly interdisciplinary topic. As a term, swarming is applied particularly to insects, but can also be applied to any other entity or animal that exhibits swarm behaviour. The term flocking or murmuration can refer specifically to swarm behaviour in birds, herding to refer to swarm behaviour in tetrapods, and shoaling or schooling to refer to swarm behaviour in fish. Phytoplankton also gather in huge swarms called blooms, although these organisms are algae and are not self-propelled the way animals are. By extension, the term "swarm" is applied also to inanimate entities which exhibit parallel behaviours, as in a robot swarm, an earthquake swarm, or a swarm of stars. From a more abstract point of view, swarm behaviour is the collective motion of a large number of self-propelled entities. From the perspective of the mathematical modeller, it is an emergent behaviour arising from simple rules that are followed by individuals and does not involve any central coordination. Swarm behaviour is also studied by active matter physicists as a phenomenon which is not in thermodynamic equilibrium, and as such requires the development of tools beyond those available from the statistical physics of systems in thermodynamic equilibrium. In this regard, swarming has been compared to the mathematics of superfluids, specifically in the context of starling flocks (murmuration). Swarm behaviour was first simulated on a computer in 1986 with the simulation program boids. This program simulates simple agents (boids) that are allowed to move according to a set of basic rules. The model was originally designed to mimic the flocking behaviour of birds, but it can be applied also to schooling fish and other swarming entities. Models In recent decades, scientists have turned to modeling swarm behaviour to gain a deeper understanding of the behaviour. Mathematical models Early studies of swarm behaviour employed mathematical models to simulate and understand the behaviour. The simplest mathematical models of animal swarms generally represent individual animals as following three rules: Move in the same direction as their neighbours Remain close to their neighbours Avoid collisions with their neighbours The boids computer program, created by Craig Reynolds in 1986, simulates swarm behaviour following the above rules. Many subsequent and current models use variations on these rules, often implementing them by means of concentric "zones" around each animal. In the "zone of repulsion", very close to the animal, the focal animal will seek to distance itself from its neighbours to avoid collision. Slightly further away, in the "zone of alignment", the focal animal will seek to align its direction of motion with its neighbours. In the outermost "zone of attraction", which extends as far away from the focal animal as it is able to sense, the focal animal will seek to move towards a neighbour. The shape of these zones will necessarily be affected by the sensory capabilities of a given animal. For example, the visual field of a bird does not extend behind its body. Fish rely on both vision and on hydrodynamic perceptions relayed through their lateral lines, while Antarctic krill rely both on vision and hydrodynamic signals relayed through antennae. However recent studies of starling flocks have shown that each bird modifies its position, relative to the six or seven animals directly surrounding it, no matter how close or how far away those animals are. Interactions between flocking starlings are thus based on a topological, rather than a metric, rule. It remains to be seen whether this applies to other animals. Another recent study, based on an analysis of high-speed camera footage of flocks above Rome and assuming minimal behavioural rules, has convincingly simulated a number of aspects of flock behaviour. Evolutionary models In order to gain insight into why animals evolve swarming behaviours, scientists have turned to evolutionary models that simulate populations of evolving animals. Typically these studies use a genetic algorithm to simulate evolution over many generations. These studies have investigated a number of hypotheses attempting to explain why animals evolve swarming behaviours, such as the selfish herd theory the predator confusion effect, the dilution effect, and the many eyes theory. Agents Self-organization Emergence The concept of emergence—that the properties and functions found at a hierarchical level are not present and are irrelevant at the lower levels–is often a basic principle behind self-organizing systems. An example of self-organization in biology leading to emergence in the natural world occurs in ant colonies. The queen does not give direct orders and does not tell the ants what to do. Instead, each ant reacts to stimuli in the form of chemical scents from larvae, other ants, intruders, food and buildup of waste, and leaves behind a chemical trail, which, in turn, provides a stimulus to other ants. Here each ant is an autonomous unit that reacts depending only on its local environment and the genetically encoded rules for its variety. Despite the lack of centralized decision making, ant colonies exhibit complex behaviours and have even been able to demonstrate the ability to solve geometric problems. For example, colonies routinely find the maximum distance from all colony entrances to dispose of dead bodies. Stigmergy A further key concept in the field of swarm intelligence is stigmergy. Stigmergy is a mechanism of indirect coordination between agents or actions. The principle is that the trace left in the environment by an action stimulates the performance of a next action, by the same or a different agent. In that way, subsequent actions tend to reinforce and build on each other, leading to the spontaneous emergence of coherent, apparently systematic activity. Stigmergy is a form of self-organization. It produces complex, seemingly intelligent structures, without need for any planning, control, or even direct communication between the agents. As such it supports efficient collaboration between extremely simple agents, who lack any memory, intelligence or even awareness of each other. Swarm intelligence Swarm intelligence is the collective behaviour of decentralized, self-organized systems, natural or artificial. The concept is employed in work on artificial intelligence. The expression was introduced by Gerardo Beni and Jing Wang in 1989, in the context of cellular robotic systems. Swarm intelligence systems are typically made up of a population of simple agents such as boids interacting locally with one another and with their environment. The agents follow very simple rules, and although there is no centralized control structure dictating how individual agents should behave, local, and to a certain degree random, interactions between such agents lead to the emergence of intelligent global behaviour, unknown to the individual agents. Swarm intelligence research is multidisciplinary. It can be divided into natural swarm research studying biological systems and artificial swarm research studying human artefacts. There is also a scientific stream attempting to model the swarm systems themselves and understand their underlying mechanisms, and an engineering stream focused on applying the insights developed by the scientific stream to solve practical problems in other areas. Algorithms Swarm algorithms follow a Lagrangian approach or an Eulerian approach. The Eulerian approach views the swarm as a field, working with the density of the swarm and deriving mean field properties. It is a hydrodynamic approach, and can be useful for modelling the overall dynamics of large swarms. However, most models work with the Lagrangian approach, which is an agent-based model following the individual agents (points or particles) that make up the swarm. Individual particle models can follow information on heading and spacing that is lost in the Eulerian approach. Ant colony optimization Ant colony optimization is a widely used algorithm which was inspired by the behaviours of ants, and has been effective solving discrete optimization problems related to swarming. The algorithm was initially proposed by Marco Dorigo in 1992, and has since been diversified to solve a wider class of numerical problems. Species that have multiple queens may have a queen leaving the nest along with some workers to found a colony at a new site, a process akin to swarming in honeybees. Ants are behaviourally unsophisticated; collectively they perform complex tasks. Ants have highly developed sophisticated sign-based communication. Ants communicate using pheromones; trails are laid that can be followed by other ants. Routing problem ants drop different pheromones used to compute the "shortest" path from source to destination(s). Self-propelled particles The concept of self-propelled particles (SPP) was introduced in 1995 by Tamás Vicsek et al. as a special case of the boids model introduced in 1986 by Reynolds. An SPP swarm is modelled by a collection of particles that move with a constant speed and respond to random perturbations by adopting at each time increment the average direction of motion of the other particles in their local neighbourhood. Simulations demonstrate that a suitable "nearest neighbour rule" eventually results in all the particles swarming together, or moving in the same direction. This emerges, even though there is no centralized coordination, and even though the neighbours for each particle constantly change over time. SPP models predict that swarming animals share certain properties at the group level, regardless of the type of animals in the swarm. Swarming systems give rise to emergent behaviours which occur at many different scales, some of which are both universal and robust. It has become a challenge in theoretical physics to find minimal statistical models that capture these behaviours. Particle swarm optimization Particle swarm optimization is another algorithm widely used to solve problems related to swarms. It was developed in 1995 by Kennedy and Eberhart and was first aimed at simulating the social behaviour and choreography of bird flocks and fish schools. The algorithm was simplified and it was observed to be performing optimization. The system initially seeds a population with random solutions. It then searches in the problem space through successive generations using stochastic optimization to find the best solutions. The solutions it finds are called particles. Each particle stores its position as well as the best solution it has achieved so far. The particle swarm optimizer tracks the best local value obtained so far by any particle in the local neighbourhood. The remaining particles then move through the problem space following the lead of the optimum particles. At each time iteration, the particle swarm optimiser accelerates each particle toward its optimum locations according to simple mathematical rules. Particle swarm optimization has been applied in many areas. It has few parameters to adjust, and a version that works well for a specific applications can also work well with minor modifications across a range of related applications. A book by Kennedy and Eberhart describes some philosophical aspects of particle swarm optimization applications and swarm intelligence. An extensive survey of applications is made by Poli. Altruism Researchers in Switzerland have developed an algorithm based on Hamilton's rule of kin selection. The algorithm shows how altruism in a swarm of entities can, over time, evolve and result in more effective swarm behaviour. Biological swarming The earliest evidence of swarm behaviour in animals dates back about 480 million years. Fossils of the trilobite Ampyx priscus have been recently described as clustered in lines along the ocean floor. The animals were all mature adults, and were all facing the same direction as though they had formed a conga line or a peloton. It has been suggested they line up in this manner to migrate, much as spiny lobsters migrate in single-file queues; it has also been suggested that the formation is the precursor for mating, as with the fly Leptoconops torrens. The findings suggest animal collective behaviour has very early evolutionary origins. Examples of biological swarming are found in bird flocks, fish schools, insect swarms, bacteria swarms, molds, molecular motors, quadruped herds and people. Social insects The behaviour of social insects (insects that live in colonies, such as ants, bees, wasps and termites) has always been a source of fascination for children, naturalists and artists. Individual insects seem to do their own thing without any central control, yet the colony as a whole behaves in a highly coordinated manner. Researchers have found that cooperation at the colony level is largely self-organized. The group coordination that emerges is often just a consequence of the way individuals in the colony interact. These interactions can be remarkably simple, such as one ant merely following the trail left by another ant. Yet put together, the cumulative effect of such behaviours can solve highly complex problems, such as locating the shortest route in a network of possible paths to a food source. The organised behaviour that emerges in this way is sometimes called swarm intelligence, a form of biological emergence. Ants Individual ants do not exhibit complex behaviours, yet a colony of ants collectively achieves complex tasks such as constructing nests, taking care of their young, building bridges and foraging for food. A colony of ants can collectively select (i.e. send most workers towards) the best, or closest, food source from several in the vicinity. Such collective decisions are achieved using positive feedback mechanisms. Selection of the best food source is achieved by ants following two simple rules. First, ants which find food return to the nest depositing a pheromone chemical. More pheromone is laid for higher quality food sources. Thus, if two equidistant food sources of different qualities are found simultaneously, the pheromone trail to the better one will be stronger. Ants in the nest follow another simple rule, to favor stronger trails, on average. More ants then follow the stronger trail, so more ants arrive at the high quality food source, and a positive feedback cycle ensures, resulting in a collective decision for the best food source. If there are two paths from the ant nest to a food source, then the colony usually selects the shorter path. This is because the ants that first return to the nest from the food source are more likely to be those that took the shorter path. More ants then retrace the shorter path, reinforcing the pheromone trail. Army ants, unlike most ant species, do not construct permanent nests; an army ant colony moves almost incessantly over the time it exists, remaining in an essentially perpetual state of swarming. Several lineages have independently evolved the same basic behavioural and ecological syndrome, often referred to as "legionary behaviour", and may be an example of convergent evolution. The successful techniques used by ant colonies have been studied in computer science and robotics to produce distributed and fault-tolerant systems for solving problems. This area of biomimetics has led to studies of ant locomotion, search engines that make use of "foraging trails", fault-tolerant storage and networking algorithms. Honey bees In temperate climates, honey bees usually form swarms in late spring. A swarm typically contains about half the workers together with the old queen, while the new queen stays back with the remaining workers in the original hive. When honey bees emerge from a hive to form a swarm, they may gather on a branch of a tree or on a bush only a few meters from the hive. The bees cluster about the queen and send out 20–50 scouts to find suitable new nest locations. The scouts are the most experienced foragers in the cluster. If a scout finds a suitable location, she returns to the cluster and promotes it by dancing a version of the waggle dance. This dance conveys information about the quality, direction, and distance of the new site. The more excited she is about her findings, the more vigorously she dances. If she can convince others they may take off and check the site she found. If they approve they may promote it as well. In this decision-making process, scouts check several sites, often abandoning their own original site to promote the superior site of another scout. Several different sites may be promoted by different scouts at first. After some hours and sometimes days, a preferred location eventually emerges from this decision-making process. When all scouts agree on the final location, the whole cluster takes off and swarms to it. Sometimes, if no decision is reached, the swarm will separate, some bees going in one direction; others, going in another. This usually results in failure, with both groups dying. A new location is typically a kilometre or more from the original hive, though some species, e.g., Apis dorsata, may establish new colonies within as little as 500 meters from the natal nest. This collective decision-making process is remarkably successful in identifying the most suitable new nest site and keeping the swarm intact. A good hive site has to be large enough to accommodate the swarm (about 15 litres in volume), has to be well-protected from the elements, receive an optimal amount of sunshine, be some height above the ground, have a small entrance and be capable of resisting ant infestation - that is why tree cavities are often selected. Non-social insects Unlike social insects, swarms of non-social insects that have been studied primarily seem to function in contexts such as mating, feeding, predator avoidance, and migration. Moths Moths may exhibit synchronized mating, during which pheromones released by females initiate searching and swarming behavior in males. Males sense pheromones with sensitive antennae and may track females as far as several kilometers away. Swarm mating involves female choice and male competition. Only one male in the swarm—typically the first—will successfully copulate. Females maximize fitness benefits and minimize cost by governing the onset and magnitude of pheromone deployed. Too little pheromone will not attract a mate, too much allows less fit males to sense the signal. After copulation, females lay the eggs on a host plant. Quality of host plant may be a factor influencing the location of swarming and egg-laying. In one case, researchers observed pink-striped oakworm moths (Anisota virginiensis) swarming at a carrion site, where decomposition likely increased soil nutrient levels and host plant quality. Flies Midges, such as Tokunagayusurika akamusi, form swarms, dancing in the air. Swarming serves multiple purposes, including the facilitation of mating by attracting females to approach the swarm, a phenomenon known as lek mating. Such cloud-like swarms often form in early evening when the sun is getting low, at the tip of a bush, on a hilltop, over a pool of water, or even sometimes above a person. The forming of such swarms is not out of instinct, but an adaptive behavior – a "consensus" – between the individuals within the swarms. It is also suggested that swarming is a ritual, because there is rarely any male midge by itself and not in a swarm. This could have formed due to the benefit of lowering inbreeding by having males of various genes gathering in one spot. The genus Culicoides, also known as biting midges, have displayed swarming behavior which are believed to cause confusion in predators. Cockroaches Cockroaches leave chemical trails in their feces as well as emitting airborne pheromones for mating. Other cockroaches will follow these trails to discover sources of food and water, and also discover where other cockroaches are hiding. Thus, groups of cockroaches can exhibit emergent behaviour, in which group or swarm behaviour emerges from a simple set of individual interactions. Cockroaches are mainly nocturnal and will run away when exposed to light. A study tested the hypothesis that cockroaches use just two pieces of information to decide where to go under those conditions: how dark it is and how many other cockroaches there are. The study conducted by José Halloy and colleagues at the Free University of Brussels and other European institutions created a set of tiny robots that appear to the roaches as other roaches and can thus alter the roaches' perception of critical mass. The robots were also specially scented so that they would be accepted by the real roaches. Locusts Locusts are the swarming phase of the short-horned grasshoppers of the family Acrididae. Some species can breed rapidly under suitable conditions and subsequently become gregarious and migratory. They form bands as nymphs and swarms as adults—both of which can travel great distances, rapidly stripping fields and greatly damaging crops. The largest swarms can cover hundreds of square miles and contain billions of locusts. A locust can eat its own weight (about 2 grams) in plants every day. That means one million locusts can eat more than one tonne of food each day, and the largest swarms can consume over 100,000 tonnes each day. Swarming in locusts has been found to be associated with increased levels of serotonin which causes the locust to change colour, eat much more, become mutually attracted, and breed much more easily. Researchers propose that swarming behaviour is a response to overcrowding and studies have shown that increased tactile stimulation of the hind legs or, in some species, simply encountering other individuals causes an increase in levels of serotonin. The transformation of the locust to the swarming variety can be induced by several contacts per minute over a four-hour period. Notably, an innate predisposition to aggregate has been found in hatchlings of the desert locust, Schistocerca gregaria, independent of their parental phase. An individual locust's response to a loss of alignment in the group appears to increase the randomness of its motion, until an aligned state is again achieved. This noise-induced alignment appears to be an intrinsic characteristic of collective coherent motion. Migratory behavior Insect migration is the seasonal movement of insects, particularly those by species of dragonflies, beetles, butterflies, and moths. The distance can vary from species to species, but in most cases these movements involve large numbers of individuals. In some cases the individuals that migrate in one direction may not return and the next generation may instead migrate in the opposite direction. This is a significant difference from bird migration. Monarch butterflies are especially noted for their lengthy annual migration. In North America they make massive southward migrations starting in August until the first frost. A northward migration takes place in the spring. The monarch is the only butterfly that migrates both north and south as the birds do on a regular basis. But no single individual makes the entire round trip. Female monarchs deposit eggs for the next generation during these migrations. The length of these journeys exceeds the normal lifespan of most monarchs, which is less than two months for butterflies born in early summer. The last generation of the summer enters into a non-reproductive phase known as diapause and may live seven months or more. During diapause, butterflies fly to one of many overwintering sites. The generation that overwinters generally does not reproduce until it leaves the overwintering site sometime in February and March. It is the second, third and fourth generations that return to their northern locations in the United States and Canada in the spring. How the species manages to return to the same overwintering spots over a gap of several generations is still a subject of research; the flight patterns appear to be inherited, based on a combination of the position of the sun in the sky and a time-compensated Sun compass that depends upon a circadian clock that is based in their antennae. Birds Bird migration Approximately 1800 of the world's 10,000 bird species are long-distance migrants. The primary motivation for migration appears to be food; for example, some hummingbirds choose not to migrate if fed through the winter. Also, the longer days of the northern summer provide extended time for breeding birds to feed their young. This helps diurnal birds to produce larger clutches than related non-migratory species that remain in the tropics. As the days shorten in autumn, the birds return to warmer regions where the available food supply varies little with the season. These advantages offset the high stress, physical exertion costs, and other risks of the migration such as predation. Many birds migrate in flocks. For larger birds, it is assumed that flying in flocks reduces energy costs. The V formation is often supposed to boost the efficiency and range of flying birds, particularly over long migratory routes. All the birds except the first fly in the upwash from one of the wingtip vortices of the bird ahead. The upwash assists each bird in supporting its own weight in flight, in the same way a glider can climb or maintain height indefinitely in rising air. Geese flying in a V formation save energy by flying in the updraft of the wingtip vortex generated by the previous animal in the formation. Thus, the birds flying behind do not need to work as hard to achieve lift. Studies show that birds in a V formation place themselves roughly at the optimum distance predicted by simple aerodynamic theory. Geese in a V-formation may conserve 12–20% of the energy they would need to fly alone. Red knots and dunlins were found in radar studies to fly 5 km per hour faster in flocks than when they were flying alone. The birds flying at the tips and at the front are rotated in a timely cyclical fashion to spread flight fatigue equally among the flock members. The formation also makes communication easier and allows the birds to maintain visual contact with each other. Other animals may use similar drafting techniques when migrating. Lobsters, for example, migrate in close single-file formation "lobster trains", sometimes for hundreds of miles. The Mediterranean and other seas present a major obstacle to soaring birds, which must cross at the narrowest points. Massive numbers of large raptors and storks pass through areas such as Gibraltar, Falsterbo, and the Bosphorus at migration times. More common species, such as the European honey buzzard, can be counted in hundreds of thousands in autumn. Other barriers, such as mountain ranges, can also cause funnelling, particularly of large diurnal migrants. This is a notable factor in the Central American migratory bottleneck. This concentration of birds during migration can put species at risk. Some spectacular migrants have already gone extinct, the most notable being the passenger pigeon. During migration the flocks were a mile (1.6 km) wide and 300 miles (500 km) long, taking several days to pass and containing up to a billion birds. Marine life Fish The term "shoal" can be used to describe any group of fish, including mixed-species groups, while "school" is used for more closely knit groups of the same species swimming in a highly synchronised and polarised manner. Fish derive many benefits from shoaling behaviour including defence against predators (through better predator detection and by diluting the chance of capture), enhanced foraging success, and higher success in finding a mate. It is also likely that fish benefit from shoal membership through increased hydrodynamic efficiency. Fish use many traits to choose shoalmates. Generally they prefer larger shoals, shoalmates of their own species, shoalmates similar in size and appearance to themselves, healthy fish, and kin (when recognised). The "oddity effect" posits that any shoal member that stands out in appearance will be preferentially targeted by predators. This may explain why fish prefer to shoal with individuals that resemble them. The oddity effect would thus tend to homogenise shoals. One puzzling aspect of shoal selection is how a fish can choose to join a shoal of animals similar to themselves, given that it cannot know its own appearance. Experiments with zebrafish have shown that shoal preference is a learned ability, not innate. A zebrafish tends to associate with shoals that resemble shoals in which it was reared, a form of imprinting. Other open questions of shoaling behaviour include identifying which individuals are responsible for the direction of shoal movement. In the case of migratory movement, most members of a shoal seem to know where they are going. In the case of foraging behaviour, captive shoals of golden shiner (a kind of minnow) are led by a small number of experienced individuals who knew when and where food was available. Radakov estimated herring schools in the North Atlantic can occupy up to with fish densities between 0.5 and 1.0 fish/cubic metre, totalling several billion fish in one school. Partridge BL (1982) "The structure and function of fish schools" Scientific American, June:114–123. Fish migration Between May and July huge numbers of sardines spawn in the cool waters of the Agulhas Bank and then follow a current of cold water northward along the east coast of South Africa. This great migration, called the sardine run, creates spectacular feeding frenzies along the coastline as marine predators, such as dolphins, sharks and gannets attack the schools. Krill Most krill, small shrimp-like crustaceans, form large swarms, sometimes reaching densities of 10,000–60,000 individual animals per cubic metre. Swarming is a defensive mechanism, confusing smaller predators that would like to pick out single individuals. The largest swarms are visible from space and can be tracked by satellite. One swarm was observed to cover an area of 450 square kilometres (175 square miles) of ocean, to a depth of 200 meters (650 feet) and was estimated to contain over 2 million tons of krill. Recent research suggests that krill do not simply drift passively in these currents but actually modify them. Krill typically follow a diurnal vertical migration. By moving vertically through the ocean on a 12-hour cycle, the swarms play a major part in mixing deeper, nutrient-rich water with nutrient-poor water at the surface. Until recently it has been assumed that they spend the day at greater depths and rise during the night toward the surface. It has been found that the deeper they go, the more they reduce their activity, apparently to reduce encounters with predators and to conserve energy. Later work suggested that swimming activity in krill varied with stomach fullness. Satiated animals that had been feeding at the surface swim less actively and therefore sink below the mixed layer. As they sink they produce faeces which may mean that they have an important role to play in the Antarctic carbon cycle. Krill with empty stomachs were found to swim more actively and thus head towards the surface. This implies that vertical migration may be a bi- or tri-daily occurrence. Some species form surface swarms during the day for feeding and reproductive purposes even though such behaviour is dangerous because it makes them extremely vulnerable to predators. Dense swarms may elicit a feeding frenzy among fish, birds and mammal predators, especially near the surface. When disturbed, a swarm scatters, and some individuals have even been observed to moult instantaneously, leaving the exuvia behind as a decoy. In 2012, Gandomi and Alavi presented what appears to be a successful stochastic algorithm for modelling the behaviour of krill swarms. The algorithm is based on three main factors: " (i) movement induced by the presence of other individuals (ii) foraging activity, and (iii) random diffusion." Copepods Copepods are a group of tiny crustaceans found in the sea and lakes. Many species are planktonic (drifting in sea waters), and others are benthic (living on the ocean floor). Copepods are typically long, with a teardrop shaped body and large antennae. Although like other crustaceans they have an armoured exoskeleton, they are so small that in most species this thin armour, and the entire body, is almost totally transparent. Copepods have a compound, median single eye, usually bright red, in the centre of the transparent head. Copepods also swarm. For example, monospecific swarms have been observed regularly around coral reefs and sea grass, and in lakes. Swarms densities were about one million copepods per cubic metre. Typical swarms were one or two metres in diameter, but some exceeded 30 cubic metres. Copepods need visual contact to keep together, and they disperse at night. Spring produces blooms of swarming phytoplankton which provide food for copepods. Planktonic copepods are usually the dominant members of the zooplankton, and are in turn major food organisms for many other marine animals. In particular, copepods are prey to forage fish and jellyfish, both of which can assemble in vast, million-strong swarms. Some copepods have extremely fast escape responses when a predator is sensed and can jump with high speed over a few millimetres (see animated image below). Planktonic copepods are important to the carbon cycle. Some scientists say they form the largest animal biomass on earth. They compete for this title with Antarctic krill. Because of their smaller size and relatively faster growth rates, however, and because they are more evenly distributed throughout more of the world's oceans, copepods almost certainly contribute far more to the secondary productivity of the world's oceans, and to the global ocean carbon sink than krill, and perhaps more than all other groups of organisms together. The surface layers of the oceans are currently believed to be the world's largest carbon sink, absorbing about 2 billion tonnes of carbon a year, the equivalent to perhaps a third of human carbon emissions, thus reducing their impact. Many planktonic copepods feed near the surface at night, then sink into deeper water during the day to avoid visual predators. Their moulted exoskeletons, faecal pellets and respiration at depth all bring carbon to the deep sea. Algal blooms Many single-celled organisms called phytoplankton live in oceans and lakes. When certain conditions are present, such as high nutrient or light levels, these organisms reproduce explosively. The resulting dense swarm of phytoplankton is called an algal bloom. Blooms can cover hundreds of square kilometres and are easily seen in satellite images. Individual phytoplankton rarely live more than a few days, but blooms can last weeks. Plants Scientists have attributed swarm behavior to plants for hundreds of years. In his 1800 book, Phytologia: or, The philosophy of agriculture and gardening, Erasmus Darwin wrote that plant growth resembled swarms observed elsewhere in nature. While he was referring to more broad observations of plant morphology, and was focused on both root and shoot behavior, recent research has supported this claim. Plant roots, in particular, display observable swarm behavior, growing in patterns that exceed the statistical threshold for random probability, and indicate the presence of communication between individual root apexes. The primary function of plant roots is the uptake of soil nutrients, and it is this purpose which drives swarm behavior. Plants growing in close proximity have adapted their growth to assure optimal nutrient availability. This is accomplished by growing in a direction that optimizes the distance between nearby roots, thereby increasing their chance of exploiting untapped nutrient reserves. The action of this behavior takes two forms: maximization of distance from, and repulsion by, neighboring root apexes. The transition zone of a root tip is largely responsible for monitoring for the presence of soil-borne hormones, signaling responsive growth patterns as appropriate. Plant responses are often complex, integrating multiple inputs to inform an autonomous response. Additional inputs that inform swarm growth includes light and gravity, both of which are also monitored in the transition zone of a root's apex. These forces act to inform any number of growing "main" roots, which exhibit their own independent releases of inhibitory chemicals to establish appropriate spacing, thereby contributing to a swarm behavior pattern. Horizontal growth of roots, whether in response to high mineral content in soil or due to stolon growth, produces branched growth that establish to also form their own, independent root swarms. Bacteria Swarming also describes groupings of some kinds of predatory bacteria such as myxobacteria. Myxobacteria swarm together in "wolf packs", actively moving using a process known as bacterial gliding and keeping together with the help of intercellular molecular signals. Mammals People A collection of people can also exhibit swarm behaviour, such as pedestrians or soldiers swarming the parapets. In Cologne, Germany, two biologists from the University of Leeds demonstrated flock like behaviour in humans. The group of people exhibited similar behavioural pattern to a flock, where if five percent of the flock changed direction the others would follow. If one person was designated as a predator and everyone else was to avoid him, the flock behaved very much like a school of fish. Understanding how humans interact in crowds is important if crowd management is to effectively avoid casualties at football grounds, music concerts and subway stations. The mathematical modelling of flocking behaviour is a common technology, and has found uses in animation. Flocking simulations have been used in many films to generate crowds which move realistically. Tim Burton's Batman Returns was the first movie to make use of swarm technology for rendering, realistically depicting the movements of a group of bats using the boids system. The Lord of the Rings film trilogy made use of similar technology, known as Massive, during battle scenes. Swarm technology is particularly attractive because it is cheap, robust, and simple. An ant-based computer simulation using only six interaction rules has also been used to evaluate aircraft boarding behaviour. Airlines have also used ant-based routing in assigning aircraft arrivals to airport gates. An airline system developed by Douglas A. Lawson uses swarm theory, or swarm intelligence—the idea that a colony of ants works better than one alone. Each pilot acts like an ant searching for the best airport gate. "The pilot learns from his experience what's the best for him, and it turns out that that's the best solution for the airline," Lawson explains. As a result, the "colony" of pilots always go to gates they can arrive and depart quickly. The program can even alert a pilot of plane back-ups before they happen. "We can anticipate that it's going to happen, so we'll have a gate available," says Lawson. Swarm behaviour occurs also in traffic flow dynamics, such as the traffic wave. Bidirectional traffic can be observed in ant trails. In recent years this behaviour has been researched for insight into pedestrian and traffic models. Simulations based on pedestrian models have also been applied to crowds which stampede because of panic. Herd behaviour in marketing has been used to explain the dependencies of customers' mutual behaviour. The Economist reported a recent conference in Rome on the subject of the simulation of adaptive human behaviour. It shared mechanisms to increase impulse buying and get people "to buy more by playing on the herd instinct." The basic idea is that people will buy more of products that are seen to be popular, and several feedback mechanisms to get product popularity information to consumers are mentioned, including smart card technology and the use of Radio Frequency Identification Tag technology. A "swarm-moves" model was introduced by a Florida Institute of Technology researcher, which is appealing to supermarkets because it can "increase sales without the need to give people discounts." Robotics The application of swarm principles to robots is called swarm robotics, while swarm intelligence refers to the more general set of algorithms. Partially inspired by colonies of insects such as ants and bees, researchers are modelling the behaviour of swarms of thousands of tiny robots which together perform a useful task, such as finding something hidden, cleaning, or spying. Each robot is quite simple, but the emergent behaviour of the swarm is more complex. The whole set of robots can be considered as one single distributed system, in the same way an ant colony can be considered a superorganism, exhibiting swarm intelligence. The largest swarms so far created is the 1024 robot Kilobot swarm. Other large swarms include the iRobot swarm, the SRI International/ActivMedia Robotics Centibots project, and the Open-source Micro-robotic Project swarm, which are being used to research collective behaviours. Swarms are also more resistant to failure. Whereas one large robot may fail and ruin a mission, a swarm can continue even if several robots fail. This could make them attractive for space exploration missions, where failure is normally extremely costly. In addition to ground vehicles, swarm robotics includes also research of swarms of aerial robots and heterogeneous teams of ground and aerial vehicles. In contrast macroscopic robots, colloidal particles at microscale can also be adopted as agents to perform collective behaviors to conduct tasks using mechanical and physical approaches, such as reconfigurable tornado-like microswarm mimicking schooling fish, hierarchical particle species mimicking predating behavior of mammals, micro-object manipulation using a transformable microswarm. The fabrication of such colloidal particles is usually based on chemical synthesis. Military Military swarming is a behaviour where autonomous or partially autonomous units of action attack an enemy from several different directions and then regroup. Pulsing, where the units shift the point of attack, is also a part of military swarming. Military swarming involves the use of a decentralized force against an opponent, in a manner that emphasizes mobility, communication, unit autonomy and coordination or synchronization. Historically military forces used principles of swarming without really examining them explicitly, but now active research consciously examines military doctrines that draw ideas from swarming. Merely because multiple units converge on a target, they are not necessarily swarming. Siege operations do not involve swarming, because there is no manoeuvre; there is convergence but on the besieged fortification. Nor do guerrilla ambushes constitute swarms, because they are "hit-and-run". Even though the ambush may have several points of attack on the enemy, the guerillas withdraw when they either have inflicted adequate damage, or when they are endangered. In 2014 the U. S. Office of Naval Research released a video showing tests of a swarm of small autonomous drone attack boats that can steer and take coordinated offensive action as a group. Gallery Myths There is a popular myth that lemmings commit mass suicide by swarming off cliffs when they migrate. Driven by strong biological urges, some species of lemmings may migrate in large groups when population density becomes too great. Lemmings can swim and may choose to cross a body of water in search of a new habitat. In such cases, many may drown if the body of water is so wide as to stretch their physical capability to the limit. This fact combined with some unexplained fluctuations in the population of Norwegian lemmings gave rise to the myth. Piranha have a reputation as fearless fish that swarm in ferocious and predatory packs. However, recent research, which started "with the premise that they school as a means of cooperative hunting", discovered that they were in fact rather fearful fish, like other fish, who schooled for protection from their predators, such as cormorants, caimans and dolphins. A researcher described them as "basically like regular fish with large teeth".	Biology and health sciences	Ethology	Biology
207888	https://en.wikipedia.org/wiki/Atlantic%20herring	Atlantic herring	Atlantic herring (Clupea harengus) is a herring in the family Clupeidae. It is one of the most abundant fish species in the world. Atlantic herrings can be found on both sides of the Atlantic Ocean, congregating in large schools. They can grow up to in length and weigh up to . They feed on copepods, krill and small fish, while their natural predators are seals, whales, cod and other larger fish. The Atlantic herring fishery has long been an important part of the economy of New England and the Atlantic provinces of Canada. This is because the fish congregate relatively near to the coast in massive schools, notably in the cold waters of the semi-enclosed Gulf of Maine and Gulf of St. Lawrence. North Atlantic herring schools have been measured up to in size, containing an estimated four billion fish. Description Atlantic herring have a fusiform body. Gill rakers in their mouths filter incoming water, trapping any zooplankton and phytoplankton. Atlantic herring are in general fragile. They have large and delicate gill surfaces, and contact with foreign matter can strip away their large scales. They have retreated from many estuaries worldwide due to excess water pollution although in some estuaries that have been cleaned up, herring have returned. The presence of their larvae indicates cleaner and more–oxygenated waters. Range and habitat Atlantic herring can be found on both sides of the Atlantic Ocean. They range, shoaling and schooling, across North Atlantic waters such as the Gulf of Maine, the Gulf of St Lawrence, the Bay of Fundy, the Labrador Sea, the Davis Straits, the Beaufort Sea, the Denmark Strait, the Norwegian Sea, the North Sea, the Skagerrak, the English Channel, the Celtic Sea, the Irish Sea, the Bay of Biscay and Sea of the Hebrides. Although Atlantic herring are found in the northern waters surrounding the Arctic, they are not considered to be an Arctic species. Baltic herring The small-sized herring in the inner parts of the Baltic Sea, which is also less fatty than the true Atlantic herring (Clupea harengus harengus), is considered a distinct subspecies, "Baltic herring" (Clupea harengus membras), despite the lack of a distinctive genome. The Baltic herring has a specific name in many local languages (Swedish strömming, Finnish silakka, Estonian räim, silk, Livonian siļk, Russian салака, Polish śledź bałtycki, Latvian reņģes, Lithuanian strimelė) and is popularly and in cuisine considered distinct from herring. For example, the Swedish dish surströmming is made from Baltic herring. Fisheries for Baltic herring have been at unsustainable levels since the Middle Ages. Around this time, the primary Baltic herring catch consisted of an autumn-spawning population. Cooling in the mid-16th century related to the Little Ice Age, combined with this overfishing, led to a dramatic loss of productivity in the population of autumn-spawning herring that rendered it nearly extinct. Due to this, the autumn-spawning herring were largely replaced by a spring-spawning population, which has since comprised most of the Baltic herring fisheries; this population is also at risk of overfishing. Life cycle Herrings reach sexual maturity when they are 3 to 5 years old. The life expectancy once mature is 12 to 16 years. Atlantic herring may have different spawning components within a single stock which spawn during different seasons. They spawn in estuaries, coastal waters or in offshore banks. Fertilization is external, as in most other fish: the female releases between 20,000 and 40,000 eggs and the males simultaneously release masses of milt so that they mix freely in the sea. Once fertilized the 1 to 1.4 mm diameter eggs sink to the sea bed where their sticky surface adheres to gravel or weed. They mature in 1–3 weeks; in 14–19 °C water it takes 6–8 days, in 7,5 °C it takes 17 days. They will only mature if the water temperature stays below 19 °C. The hatched larvae are 3 to 4 mm long and transparent except for the eyes which have some pigmentation. Population Herrings are most seen in the North Atlantic Ocean, from the coast of South Carolina until Greenland, and from the Baltic Sea until Novaya Zemlya. In the North Sea people can distinguish four different main populations spawning in different periods: The Buchan-Shetland herrings spawn in August and September near the Scottish and Shetland coasts. On the Dogger Bank herrings spawn from August until October. The more southern population will spawn later, from November until January. These are the herrings from the Southern Bight of Downs. The Soused herring spawns every spring in the Baltic Sea, and travels via Skagerrak to the North Sea. These four populations live outside of the spawn season interchangeably. In their spawn season, each population gathers together on their own spawn grounds. In the past, there was another, fifth distinct population, the Zuiderzee herring, which spawned in the former Zuiderzee. This population disappeared when the Zuiderzee was drained by the Dutch as part of the larger Zuiderzee Works. Ecology Herring-like fish are the most important fish group on the planet. They are also the most populous fish. They are the dominant converter of zooplankton into fish, consuming copepods, arrow worms chaetognatha, pelagic amphipods hyperiidae, mysids and krill in the pelagic zone. Conversely, they are a central prey item or forage fish for higher trophic levels. The reasons for this success are still enigmatic; one speculation attributes their dominance to the huge, extremely fast cruising schools they inhabit. Orca, cod, dolphins, porpoises, sharks, rockfish, seabirds, whales, squid, sea lions, seals, tuna, salmon, and fishermen are among the predators of these fishes. Herring's pelagic–prey includes copepods (e.g. Centropagidae, Calanus spp., Acartia spp., Temora spp.), amphipods like Hyperia spp., larval snails, diatoms by larvae below , peridinians, molluscan larvae, fish eggs, krill like Meganyctiphanes norvegica, mysids, small fishes, menhaden larvae, pteropods, annelids, tintinnids by larvae below , Haplosphaera, Pseudocalanus. Schooling Atlantic herring can school in immense numbers. Radakov estimated herring schools in the North Atlantic can occupy up to 4.8 cubic kilometres with fish densities between 0.5 and 1.0 fish/cubic metre, equivalent to several million fish in one school. Herring are amongst the most spectacular schoolers ("obligate schoolers" under older terminology). They aggregate in groups that consist of thousands to hundreds of thousands or even millions of individuals. The schools traverse the open oceans. Schools have a very precise spatial arrangement that allows the school to maintain a relatively constant cruising speed. Schools from an individual stock generally travel in a triangular pattern between their spawning grounds, e.g. Southern Norway, their feeding grounds (Iceland) and their nursery grounds (Northern Norway). Such wide triangular journeys are probably important because feeding herrings cannot distinguish their own offspring. They have excellent hearing, and a school can react very quickly to evade predators. Herring schools keep a certain distance from a moving scuba diver or a cruising predator like a killer whale, forming a vacuole which looks like a doughnut from a spotter plane. The phenomenon of schooling is far from understood, especially the implications on swimming and feeding-energetics. Many hypotheses have been put forward to explain the function of schooling, such as predator confusion, reduced risk of being found, better orientation, and synchronized hunting. However, schooling has disadvantages such as: oxygen- and food-depletion and excretion buildup in the breathing media. The school-array probably gives advantages in energy saving although this is a highly controversial and much debated field. Schools of herring can on calm days sometimes be detected at the surface from more than a mile away by the little waves they form, or from a few meters at night when they trigger bioluminescence in surrounding plankton ("firing"). All underwater recordings show herring constantly cruising reaching speeds up to per second, and much higher escape speeds. Relationship with humans Fisheries The Atlantic herring fishery is managed by multiple organizations that work together on the rules and regulations applying to herring. As of 2010 the species was not threatened by overfishing. They are an important bait fish for recreational fishermen. Aquariums Because of their feeding habits, cruising desire, collective behavior and fragility they survive in very few aquaria worldwide despite their abundance in the ocean. Even the best facilities leave them slim and slow compared to healthy wild schools.	Biology and health sciences	Clupeiformes	Animals
207897	https://en.wikipedia.org/wiki/Window%20%28computing%29	Window (computing)	In computing, a window is a graphical control element. It consists of a visual area containing some of the graphical user interface of the program it belongs to and is framed by a window decoration. It usually has a rectangular shape that can overlap with the area of other windows. It displays the output of and may allow input to one or more processes. Windows are primarily associated with graphical displays, where they can be manipulated with a pointer by employing some kind of pointing device. Text-only displays can also support windowing, as a way to maintain multiple independent display areas, such as multiple buffers in Emacs. Text windows are usually controlled by keyboard, though some also respond to the mouse. A graphical user interface (GUI) using windows as one of its main "metaphors" is called a windowing system, whose main components are the display server and the window manager. History The idea was developed at the Stanford Research Institute (led by Douglas Engelbart). Their earliest systems supported multiple windows, but there was no obvious way to indicate boundaries between them (such as window borders, title bars, etc.). Research continued at Xerox Corporation's Palo Alto Research Center / PARC (led by Alan Kay). They used overlapping windows. During the 1980s the term "WIMP", which stands for window, icon, menu, pointer, was coined at PARC. Apple had worked with PARC briefly at that time. Apple developed an interface based on PARC's interface. It was first used on Apple's Lisa and later Macintosh computers. Microsoft was developing Office applications for the Mac at that time. Some speculate that this gave them access to Apple's OS before it was released and thus influenced the design of the windowing system in what would eventually be called Microsoft Windows. Properties Windows are two dimensional objects arranged on a plane called the desktop metaphor. In a modern full-featured windowing system they can be resized, moved, hidden, restored or closed. Windows usually include other graphical objects, possibly including a menu-bar, toolbars, controls, icons and often a working area. In the working area, the document, image, folder contents or other main object is displayed. Around the working area, within the bounding window, there may be other smaller window areas, sometimes called panes or panels, showing relevant information or options. The working area of a single document interface holds only one main object. "Child windows" in multiple document interfaces, and tabs for example in many web browsers, can make several similar documents or main objects available within a single main application window. Some windows in macOS have a feature called a drawer, which is a pane that slides out the side of the window and to show extra options. Applications that can run either under a graphical user interface or in a text user interface may use different terminology. GNU Emacs uses the term "window" to refer to an area within its display while a traditional window, such as controlled by an X11 window manager, is called a "frame". Any window can be split into the window decoration and the window's content, although some systems purposely eschew window decoration as a form of minimalism. Window decoration The window decoration is a part of a window in most windowing systems. Window decoration typically consists of a title bar, usually along the top of each window and a minimal border around the other three sides. On Microsoft Windows this is called "non-client area". In the predominant layout for modern window decorations, the top bar contains the title of that window and buttons which perform windowing-related actions such as: Close Maximize Minimize Resize Roll-up The border exists primarily to allow the user to resize the window, but also to create a visual separation between the window's contents and the rest of the desktop environment. Window decorations are considered important for the design of the look and feel of an operating system and some systems allow for customization of the colors, styles and animation effects used. Window border Window border is a window decoration component provided by some window managers, that appears around the active window. Some window managers may also display a border around background windows. Typically window borders enable the window to be resized or moved by dragging the border. Some window managers provide useless borders which are purely for decorative purposes and offer no window motion facility. These window managers do not allow windows to be resized by using a drag action on the border. Title bar The title bar is a graphical control element and part of the window decoration provided by some window managers. As a convention, it is located at the top of the window as a horizontal bar. The title bar is typically used to display the name of the application or the name of the open document, and may provide title bar buttons for minimizing, maximizing, closing or rolling up of application windows. These functions are typically placed in the top-right of the screen to allow fast and inaccurate inputs through barrier pointing. Typically title bars can be used to provide window motion enabling the window to be moved around the screen by grabbing the title bar and dragging it. Some window managers provide title bars which are purely for decorative purposes and offer no window motion facility. These window managers do not allow windows to be moved around the screen by using a drag action on the title bar. Default title-bar text often incorporates the name of the application and/or of its developer. The name of the host running the application also appears frequently. Various methods (menu-selections, escape sequences, setup parameters, command-line options – depending on the computing environment) may exist to give the end-user some control of title-bar text. Document-oriented applications like a text editor may display the filename or path of the document being edited. Most web browsers will render the contents of the HTML element title in their title bar, sometimes pre- or postfixed by the application name. Google Chrome and some versions of Mozilla Firefox place their tabs in the title bar. This makes it unnecessary to use the main window for the tabs, but usually results in the title becoming truncated. An asterisk at its beginning may be used to signify unsaved changes. The title bar often contains widgets for system commands relating to the window, such as a maximize, minimize, rollup and close buttons; and may include other content such as an application icon, a clock, etc. Title bar buttons Some window managers provide title bar buttons which provide the facility to minimize, maximize, roll-up or close application windows. Some window managers may display the title bar buttons in the task bar or task panel, rather than in the title bars. The following buttons may appear in the title bar: Close Maximize Minimize Resize Roll-up (or WindowShade) Note that a context menu may be available from some title bar buttons or by right-clicking. Title bar icon Some window managers display a small icon in the title bar that may vary according to the application on which it appears. The title bar icon may behave like a menu button, or may provide a context menu facility. macOS applications commonly have a proxy icon next to the window title that functions the same as the document's icon in the file manager. Document status icon Some window managers display an icon or symbol to indicate that the contents of the window have not been saved or confirmed in some way: macOS displays a dot in the center of its close button; RISC OS appends an asterisk to the title. Tiling window managers Some tiling window managers provide title bars which are purely for informative purposes and offer no controls or menus. These window managers do not allow windows to be moved around the screen by using a drag action on the title bar and may also serve the purpose of a status line from stacking window managers. In popular operating systems	Technology	User interface	null
207906	https://en.wikipedia.org/wiki/Tinamou	Tinamou	Tinamous () are members of the order Tinamiformes (), and family Tinamidae (), divided into two distinct subfamilies, containing 46 species found in Mexico, Central America, and South America. The word "tinamou" comes from the Galibi term for these birds, tinamu. Tinamous are the only living group of palaeognaths able to fly, and were traditionally regarded as the sister group of the flightless ratites, but recent work places them well within the ratite radiation as most closely related to the extinct moa of New Zealand, implying flightlessness emerged among ratites multiple times. Tinamous first appear in the fossil record in the Miocene epoch. They are generally sedentary, ground-dwelling and, though not flightless, when possible avoid flight in favour of hiding or running away from danger. They are found in a variety of habitats, ranging from semi-arid alpine grasslands to tropical rainforests. The two subfamilies are broadly divided by habitat, with the Nothurinae referred to as steppe or open country tinamous, and the Tinaminae known as forest tinamous. Although some species are quite common, tinamous are shy and secretive birds. They are active during the day, retiring to roosts at night. They generally have cryptic plumage, with males and females similar in appearance, though the females are usually larger. They are opportunistic and omnivorous feeders, consuming a wide variety of plant and animal food from fruits and seeds to worms, insects and small vertebrates. They will dust-bathe as well as wash themselves by standing in heavy rain. They are heard more often than seen, communicating with each other by a variety of frequently given, characteristic calls, especially during the breeding season. With occasional exceptions, a male tinamou maintains a territory and a nesting site during the breeding season which a succession of females will visit, laying their eggs in the same nest. Females will wander through several territories mating with, and laying eggs in the nests of, the resident males. Nests are always on the ground, concealed in vegetation or among rocks. Eggs are relatively large and glossy, often brightly colored when laid, and are incubated by the males for a period of two to three weeks. The chicks can run soon after hatching and are largely self-sufficient at three weeks old. Tinamous and their eggs have many natural predators, from falcons and vampire bats to jaguars. They have also been extensively hunted by humans and sometimes persecuted as agricultural pests. However, the main threat to their populations is from habitat destruction through land clearing and agricultural development. Seven species are listed as vulnerable and another seven as near-threatened. They feature in the mythology of the indigenous peoples of their range. Often translocated and easily bred in captivity, they have never been successfully domesticated. Taxonomy and systematics The tinamou family consists of 46 extant species in nine genera. The two subfamilies are the Nothurinae (also known as the Rhyncotinae), the steppe tinamous, and the Tinaminae, the forest tinamous. "Tinamidae" was defined as by Gauthier and de Queiroz (2001): "Tinamidae refers to the crown clade stemming from the most recent common ancestor of Tetrao [Tinamus] major Gmelin 1789 and all extant birds sharing a more recent ancestor with that species than with Struthio camelus Linnaeus 1758 and Vultur gryphus Linnaeus 1758." Their similarity to other ground-dwelling birds such as partridges and megapodes is a result of convergence and symplesiomorphy rather than shared evolutionary innovations. Of Gondwanan origin, tinamous are allied to the flightless ratites, together comprising the Palaeognathae ("old jaws"), while all other living birds are members of Neognathae ("new jaws"). Unlike other palaeognaths, tinamous do have a keeled sternum, but like the other palaeognaths, they have a distinctive palate. It was formerly believed that the Tinamiformes separated from the ratites early on due to their retention of a keeled sternum. The tinamous' possession of powder-down feathers and preen glands, which the other ratites lack, was another source of confusion in evaluating their taxonomy. The tinamou family has been shown to be monophyletic. Phylogenomic studies have placed it as the sister group to extant Australasian and Oceanian ratites (i.e. the cassowaries, emus, and kiwis), thus putting it well within the ratite phylogenetic tree, with the South American rheas and African ostriches as successive outgroups. Research published starting in 2010 has found that tinamous are closest to the extinct moa of New Zealand; moa are more distantly related to the geographically proximate kiwis, emus and cassowaries than had been previously supposed. These findings imply that flightlessness evolved independently multiple times in ratite evolution. Flight may have been maintained in the tinamou family due to the rhea colonizing South America before ancestral tinamous arrived. The ecological niche for large, flightless herbivores was thus already occupied, forcing tinamous to retain smaller-bodied, omnivorous, and volant lifestyles. Fossil record Flight-capable lithornithids from the Paleocene and Eocene epochs appear to have been structurally the most similar precursors to the tinamous, and may have been ancestral to them as well as to the ratites, though their precise relationships are unclear. The earliest unequivocally Tinamiforme fossil material dates from the Miocene, but flightless ratite-like taxa from the Paleocene may belong to this group. Several tinamou fossils have been found in the 16–17 Mya Early-Middle Miocene Santa Cruz Formation and the contemporary, or slightly older, Pinturas Formation, in Santa Cruz Province of Argentinian Patagonia, including a tinaminid, Crypturellus reai. Associated fossils indicate that the local palaeoenvironment at the beginning of this period was characterised by a humid, subtropical climate, with forest vegetation, becoming drier and more open with time. Some of the tinamou fossil material appears to be intermediate between the two subfamilies, suggesting that the period coincides with the origins of the radiation of the Nothurinae into the expanding open-country habitats. Nothurine fossils referrable to Eudromia and Nothura have been found in the Late Miocene Cerro Azul Formation from the Pampean region of central-southern Argentina. Tinamous described from Pliocene material include Eudromia olsoni Tambussi & Tonni, 1985, Nothura parvulus Rovereto, 1914, and Nothura padulosa Mercerat, 1897. The Pliocene fossil genera Cayetornis Brodkorb and Tinamisornis Rovereto have been synonymized with Nothura and Eudromia respectively. Fossils having affinities with several extant genera have been found in Pleistocene deposits. Generic relationships Cladogram of tinamou genera based on a study by Lukas Musher and collaborators published in 2024. Species in taxonomic order Conservation status key: VU – vulnerable NT – near threatened LC – least concern Order Tinamiformes Huxley 1872 [Crypturi Goodchild 1891; Dromaeomorphae Huxley 1867] Family Tinamidae Genus †Querandiornis Rusconi 1958 †Querandiornis romani Rusconi 1958 Subfamily Tinaminae Genus Crypturellus †Crypturellus reai Chandler 2012 Barred tinamou, Crypturellus casiquiare – LC Bartlett's tinamou, Crypturellus bartletti – LC Berlepsch's tinamou, Crypturellus berlepschi – LC Black-capped tinamou, Crypturellus atrocapillus – NT Brazilian tinamou, Crypturellus strigulosus – LC Brown tinamou, Crypturellus obsoletus – LC Choco tinamou, Crypturellus kerriae – VU Cinereous tinamou, Crypturellus cinereus – LC Grey-legged tinamou, Crypturellus duidae – NT Little tinamou, Crypturellus soui – LC Pale-browed tinamou, Crypturellus transfasciatus – NT Red-legged tinamou, Crypturellus erythropus – LC Colombian tinamou, C. (e.) columbianus (taxonomic status presently unclear) SACC in 2006 did not approve the split, BLI followed suit. Magdalena tinamou, C. (e.) saltuarius (taxonomic status presently unclear) SACC in 2006 did not approve the split, BLI followed suit. Santa Marta tinamou, C. (e.) idoneus (taxonomic status presently unclear) SACC in 2006 did not approve the split, BLI followed suit. Rusty tinamou, Crypturellus brevirostris, also known as short-billed tinamou – LC Slaty-breasted tinamou, Crypturellus boucardi, also known as Boucard's tinamou – LC Small-billed tinamou, Crypturellus parvirostris – LC Tataupa tinamou, Crypturellus tataupa – LC Tepui tinamou, Crypturellus ptaritepui – LC Thicket tinamou, Crypturellus cinnamomeus – LC Undulated tinamou, Crypturellus undulatus – LC Variegated tinamou, Crypturellus variegatus – LC Yellow-legged tinamou, Crypturellus noctivagus – NT Genus Nothocercus Highland tinamou, Nothocercus bonapartei – LC Hooded tinamou, Nothocercus nigrocapillus – VU Tawny-breasted tinamou, Nothocercus julius – LC Genus Tinamus Black tinamou, Tinamus osgoodi – VU Great tinamou, Tinamus major – NT Grey tinamou, Tinamus tao – VU Solitary tinamou, Tinamus solitarius – NT White-throated tinamou, Tinamus guttatus – NT Subfamily Nothurinae Genus Eudromia †Eudromia intermedia (Rovereto 1914) [Tinamisornis intermedius Rovereto 1914 non Dabbene & Lillo 1913; Roveretornis intermedius (Rovereto 1914) Brodkorb 1961] †Eudromia olsoni Tambussi & Tonni 1985 [Tinamisornis intermedius Dabbene & Lillo 1913 non Rovereto 1914; Eudromia elegans intermedia (Dabbene & Lillo 1913)] Elegant crested tinamou, Eudromia elegans – LC Quebracho crested tinamou, Eudromia formosa – LC Genus Nothoprocta Andean tinamou, Nothoprocta pentlandii – LC Brushland tinamou, Nothoprocta cinerascens – LC Chilean tinamou, Nothoprocta perdicaria – LC Curve-billed tinamou, Nothoprocta curvirostris – LC Ornate tinamou, Nothoprocta ornata – LC Taczanowski's tinamou, Nothoprocta taczanowskii – VU Genus Nothura Chaco nothura, Nothura chacoensis – LC Darwin's nothura, Nothura darwinii – LC Lesser nothura, Nothura minor – VU †Nothura paludosa Mercerat 1897 †Nothura parvula (Rovereto) Tambussi 1989 [Tinamisornis parvulus Rovereto; Cayetanornis parvulus (Rovereto) Brodkorb 1963] Spotted nothura, Nothura maculosa – LC White-bellied nothura, Nothura boraquira – LC Genus Rhynchotus Huayco tinamou, Rhynchotus maculicollis – LC Red-winged tinamou, Rhynchotus rufescens – LC Genus Taoniscus Dwarf tinamou, Taoniscus nanus, also known as least tinamou – EN Genus Tinamotis Patagonian tinamou, Tinamotis ingoufi, also known as Ingouf's tinamou – LC Puna tinamou, Tinamotis pentlandii, also known as Pentland's tinamou – LC Description Tinamous are plump, compact birds with slender necks, small heads and, usually, short, decurved bills, though a few have long bills. Females are usually larger than the males. The smallest species, the dwarf tinamou, weighs about with a length of . Females of the largest, the grey tinamou, weigh up to with a length of up to . Their feet have three forward-facing toes; a hind toe is either higher and retrogressed, or absent. The back of the tarsus is covered with scales, the color of which may aid in identification. Tinamous have a pneumaticized skeleton with a sternal keel, 16–18 cervical vertebrae, and fused thoracic vertebrae. They have poor circulation, evidenced by a greenish tint to the skin. They also have relatively the smallest hearts and lungs of all birds, comprising only 1.6–3.1% of their body weight, whereas the equivalent in a domestic chicken is 12%. Despite their poor flying ability, the percentage of their body mass that is muscle is 28.6–40%, which is similar to that of hummingbirds. The preen gland is small and tufted. The male has a corkscrew shaped penis, similar to those of the other ratites and to the hemipenis of some reptiles. The female has a small phallic organ in the cloaca which becomes larger during the breeding season. Plumage The plumage of the family is cryptic, as is usual with ground birds, with typical colors ranging through dark brown, rufous, buff, yellow and grey. Plumage does not usually differ between sexes, but in a few species females are brighter. The forest dwellers tend to be darker and more uniform, whereas the steppe species are paler with more barring, speckling, or streaking. Tinamous have well-developed powder down feathers; these grow continuously and disintegrate at the tips into a powder that is spread through the rest of the feathers by preening. This gives the plumage a glossy appearance as well as waterproofing it. Their tails are short, sometimes hidden behind the coverts, and possibly indicative of an ability to sacrifice feathers to a predator in order to escape when grasped. Some tinamous have crests. Members of Eudromia have the most developed crests and, when excited, will direct them forward. Voice Tinamous are rarely seen but often heard within their range and have a wide variety of calls. They are among the most characteristic bird vocalizations of South America and Central America, often resembling sounds made by a flute or a whistle. Some calls are uniform and monotone, while others have multiple phrases. They vary in intensity and can often be heard from afar. Trying to locate a bird by its call is not easy. Plains-dwelling tinamous have higher-pitched, more delicate voices. They can also be less melodic, sometimes resembling the chirps of crickets. Forest species tend to have deep, loud calls, suitable for penetrating the vegetation. The male highland tinamou can be heard several kilometres distant through dense forest. When calling, a tinamou extends its neck vertically, tilts its head at an angle, and opens its bill wide. A bird, when flushed, will utter a sharp trill. Identification of tinamous is not an easy task; utilizing their calls as a tool is integral. Each species has its own unique call or calls. The solitary tinamou has 11 different vocalizations. In most species both sexes call; some have different calls for males and females. Females tend to have deeper voices. Some species, in particular members of Crypturellus, have regional dialects. Male slaty-breasted tinamous have calls unique enough to be individually recognized by humans. Calls are typically heard more frequently during the breeding season. However, the time of day can differ amongst species, as some are more vocal in the morning, others in the evening, and some are more vocal during the heat of midday. Some will call at night from their roosts. Frequency can vary between species and between individuals. One male brushland tinamou called every few minutes from dawn until dusk (over 500 calls daily). Some, in particular Crypturellus species, use regular call sites. Only a few possess an alarm call. Distribution and habitat Range Tinamous are exclusively neotropical and all 47 species live in South America, Mexico, and Central America. The range of the northernmost species extends to Mexico but not much further north than the Tropic of Cancer. Chilean tinamous have been introduced to Easter Island. The greatest concentration of species is in the tropics, and in particular the Amazon Basin. In the north, they tend to be forest or woodland birds, while in the south they prefer open habitats. Tinamous form the dominant group of terrestrial birds in South America, where they largely replace the Galliformes ecologically, with no other bird family there having comparable diversity, distribution, or suite of habitat adaptations. Rheas are only found in open country, curassows and guans are generally limited to forests, and the pheasant family is only represented by a few species in the north of the region. They occur in a wide range of habitats. Members of the genera Tinamus, Nothocercus, and Crypturellus live in dense forests, with Nothocercus preferring high altitude, and members of most other genera in grassland, puna, montane forest, and savanna. Tinamotis and Nothoprocta prefer high altitude habitats, up to , whereas the other steppe tinamous have a wide altitude range. Tinamous inhabit most parts of South and Central America, as well as the tropical regions of Mexico, with the exception of aquatic, snow-covered, and true desert habitats, and the southernmost tip of Patagonia. Ecology Behavioral and ecological separation of tinamou species is evident where their ranges overlap through the utilization of different food sources and occupation of limited micro-habitats. These micro-habitats are not always easy to identify, and are highly vulnerable to environmental changes. Some species, such as the red-winged tinamou, utilize multiple habitats such as the open savannas of Amazonia and the dry valleys of the Andes. Similarly, brown tinamous occur in both the Amazon basin and the humid montane forests on the Andean slope. Panama provides examples of ecological separation. The highland tinamou occupies the highlands throughout the country. The great tinamou prefers the rainforests on the slopes. The Choco tinamou also likes the rainforest, but is limited to the south-east of the country. Finally, the little tinamou is found in dense secondary forest on either the Pacific or Atlantic slope above . Size difference allows the red-winged tinamou and the spotted nothura to coexist, as they both occupy the same habitat of Brazil, the tropical savanna. The former prefers long grass pastures, while the latter prefers short grass. Further examples of such diversity are found in the Andes, where a small subspecies of Darwin's nothura, Nothura darwinii boliviana, occurs in grassland at about above sea level. Here also are the red-winged tinamou which prefers open ground with some scrub, and the Andean tinamou which prefers dense vegetation beside streams. Their habitat extends upslope through the Polylepis woodlands into puna grassland. In the puna is another subspecies of Darwin's nothura, Nothura darwinii agassizii, which prefers tussock grassland. Also in the puna is the ornate tinamou which frequents the rocky slopes and cliffs of tola heath. Higher in the puna is the puna tinamou, living just below the snowline at as well as in the semi-deserts of the southern Altiplano. Movements Tinamous are largely sedentary birds. Forest-dwelling tinamous will move short distances if climatic conditions, such as intense rain, flooding or drought force them to. Most Amazonian species will move between the varzea forests and dry land depending on water levels. The puna tinamou occupies high ridges in the Andes but, in bad weather, will move down to the valley floors. Forest species, such as the slaty-breasted tinamou, maintain large home ranges through which they move in apparently random patterns. The male brushland tinamou maintains a home territory of , but will occasionally wander outside it into those of his neighbors. Females will wander throughout multiple males' territories. The ornate tinamou lives mainly upslope in hilly puna grassland but will move each morning to the bottom of the slopes to feed and drink. Granivorous species will move daily into grain fields with some, such as Darwin's nothura, remaining in the fields until there is no food left. Open country and southern species maintain territories only during the breeding season and at other times seem to wander at random. Behavior Tinamous form one of the most terrestrial groups of flying birds, spending virtually all of their time on the ground. They walk silently, pausing frequently in mid-stride. When a potential threat is detected, a tinamou will typically freeze in one of two positions, either crouched or with its neck extended upwards. As far as possible, they will avoid resorting to flight by stealthy walking or running away from danger as well as by concealment in dense vegetation. They may then pause to observe the cause of their alarm from cover. They also hide in burrows. Their cryptic behavior has allowed them to survive or even thrive in areas where guans have been extirpated. Flight Unlike the related ratites, tinamous can fly, though poorly and reluctantly, preferring to walk or run. When forced to take to the air, they do so only for short distances at high speed. Their small wings give them a high wing loading. They take off with rapid and noisy wing beats, until they have gained sufficient altitude, then glide while slipping sideways, with an occasional further burst of flapping. Due to their near lack of a tail to serve as rudder or counterweight, tinamous are notoriously poor at steering. They regularly crash into objects on attempting to take off, sometimes with fatal consequences. They rarely fly more than and typically do so downslope where the terrain allows. They land in an upright position with upstretched neck. Some species will land running. The brushland tinamou will perform a sharp 90° turn immediately before touching down. Roosting Although tinamous are diurnal, many species become less active in the middle of the day. They rest or feed during this period, while during the night they will cease all activity. They are wary of the dark; they roost at night and have been known to roost during solar eclipses. Roosting of the larger forest species, such as those in Tinamus, occurs in trees. They prefer horizontal branches approximately off the ground, choosing sites with good views and clear exits. In order to minimize the effort involved in ascending to their roosts, in hilly terrain they will access them from uphill and, when threatened, will fly downhill to gain more distance from the threat. Tinamous prefer thick branches on which to roost as they do not clutch the branch with their toes, but rest on it with folded legs. They will reuse the same locations and avoid defecating nearby to avoid advertising the roost site to predators. The smaller forest species, along with the steppe tinamous, will roost on the ground, sometimes in the shelter of a bush. They will also use the same location repeatedly; known examples are the elegant crested and ornate tinamous. Sociality Tinamous, depending on the species, may be solitary or social and gather in groups. Gregariousness also varies by season. Forest species tend to be solitary and may only approach other birds during the breeding season. Some live as mated pairs throughout the year. Steppe or grassland species tend to live in groups, though with little obvious group interaction apart from an occasional contact call. Group size may vary by season; in winter, aggregations of elegant crested tinamous may approach 100 birds. Both steppe and forest species are territorial, though territoriality varies between species from being characteristic only during the breeding season, to being territorial throughout the year. When defending their territories from conspecifics, tinamous are highly vocal, creating a cacophony of sound. When an intruder is noticed, birds of the same sex will confront it. This may lead to conflict, with feet and wings being used in attack. Both males and females will defend their territories; however, in each species only one sex is fiercely territorial. Breeding In most tinamou species, the males practice simultaneous polygyny and the females sequential polyandry. This is not invariable; ornate tinamous form stable pairs, and spotted nothuras are monogamous when young and polygamous when older. There are larger numbers of females than males; for example, the variegated tinamou has a female to male ratio of 4:1. The breeding season varies from species to species; those that live in tropical forests, where there is little seasonal change, may breed at any time, though there is usually a preferred period. In areas with a marked seasonal fluctuation, tinamous generally breed when food is most abundant, which is usually summer. Studies have shown that it is not day length that determines the onset of breeding, but the amount of light, through cloud cover. Courtship The courtship process starts with the male vocally advertising his abilities with continuous calling. He will try to attract multiple females. In Tinamus species, the male will lower his chest to the ground, stretch his neck forward, and fluff up his back to appear larger than normal. When observed head on, all of the bird's back is in view while the under-tail coverts are exposed, a pose similar to that used by the rhea. The female will scratch her feet on the ground as part of the ritual. Nesting Tinamous always nest on the ground; in open areas, near a bush; in scrub, in a dense patch of grass; in forest, at the base of a tree trunk between the buttresses. The highland tinamou is unique in that it sites its nest in a cavity or under an overhanging rock on a steep slope. Many species do not build a nest, choosing to lay their eggs on a thin bed of leaves. Other species do construct nests and are meticulous in doing so. The nest of the ornate tinamou is circular and made of grass on a turf surface. The male brushland tinamou starts to scrape out a nest once copulation has occurred; several may be constructed though only one is used. Egg-laying A tinamou female lays several eggs which the male incubates while the female departs to seek another mate. Large species will lay one egg every 3–4 days, while the smaller ones lay on consecutive days. The females lay eggs in multiple nests throughout the nesting season. There may be as many as 16 eggs in a clutch, a consequence of several females laying in the same nest. The more mature male will attract more females and may have the eggs of up to four females under him. The variegated and ornate tinamous have single-female nests, and consequently only one or two eggs per nest. This may result from food shortage in their ranges and the consequent ability to care for only one or two chicks. The eggs are fairly deeply colored, usually in a single color, and have a hard porcelain-like gloss. Colors vary with species, ranging through green, purple, violet, turquoise, steel grey, chocolate and lemon-yellow. White is rare, but does occur. Though the eggs are bright and colorful when laid, over time they fade and become duller. For example, the egg of the red-winged tinamou dulls from purple to leaden. Most tinamou eggs are solid colored, without spots or speckling; however, the eggs of Tinamotis species may exhibit small white speckles. The benefit of laying brightly colored eggs is unknown, but is not detrimental as most tinamou predators hunt at night. Eggs are relatively large compared to the mass of the female, though even the largest birds produce eggs very similar in size to the smallest of species. Their shapes are either spherical or elliptical; the two ends are similar in shape, and difficult to distinguish. The shells are thin enough to see the embryos within. Incubation Incubation takes about 16 days in Crypturellus, which contains the smallest species, and 19–20 days in Tinamus and Eudromia. During this period the male is typically silent; if he does call, he does so away from the nest. As he incubates, he will leave the nest to feed, and he may be gone from 45 minutes to five hours, covering the eggs when he leaves. While incubating, he is mainly motionless and reluctant to move, even from potential danger. It is possible for a human observer to approach and touch the incubating male without eliciting an overt response. Some species will flatten themselves against the ground, stretch out their necks, and raise their backs to the air. This posture causes them to resemble a plant; however, if it is overdone, the eggs become visible from behind. If the male becomes alarmed enough to leave the nest, he will attempt a distraction display. This usually involves a fake injury display, similar to that of the killdeer. To do this, he will hop on one leg and attempt to fly, always falling down. He will perform this display if the eggs have not hatched or the chicks are still too young to fly. It is generally believed that tinamous are not as effective at distraction displays as other birds. Chicks Chicks hatch synchronously with a dense downy coat. The coloring is white, grey or yellow, with dark spots to aid in camouflage. The young are precocial, and can run almost as soon as they hatch. Soon after hatching the eggs, the male will leave the nest and call the chicks to him with a soft contact call. If threatened, he will freeze and attempt to hide the chicks under his wings or belly. There have been documented cases of females caring for the young; it is thought that this occurs when the male has been killed. Young chicks can feed themselves within the first few days, but the male will bring the food and drop it on the ground in front of them. The chicks have a high initial mortality rate. However, within a few days they are chasing insects on their own and, at 1–3 weeks, they can fly to branches a metre from the ground. They are self-sufficient within 20 days. By 20 days, the young slaty-breasted tinamou has gained adult size, though not adult weight. The spotted nothura will go from 10% of adult weight to 90% within 85 days, and the red-winged tinamou will do so in 108 days. Sexual maturity comes at the age of one year, although some species may be physiologically mature by 57 days. However, some behavior may need to be learned before the birds can breed successfully. Once done with the brood the male, if still within the breeding season, will seek out another female and initiate the cycle again. Studies have shown that 54–62% of breeding female spotted nothura are first-year birds. Feeding Foods Tinamous are opportunistic feeders and eat a wide range of foods, though each species varies in proportional dietary makeup. Tinamou genera can be roughly divided into three groups based on the vegetable component of their diets. Tinamus, Nothocercus and Crypturellus focus on fleshy fruit. Nothura, Nothoprocta and Eudromia, comprising open country birds, eat mainly seeds and other soft vegetative matter. High-altitude genera living in harsh environments, such as Tinamotis, will eat most of the plant, not just the succulent parts. Most species eat a mixture of plant and animal products, though some are mainly herbivorous and others predominantly insectivorous or carnivorous. Diet may also vary seasonally; red-winged tinamous eat mostly animal food in the summer and plant matter in the winter. Chicks eat more insects than their parents, probably for their growth needs. Consumed plant material includes fruit (either fallen or on the tree), seeds, green shoots, tender leaves, buds, flowers, tender stems, roots, and tubers. Much of the animal food consists of insects, including ants, termites, beetles, grasshoppers, Hemiptera, and lepidopteran larvae, as well as gastropods, mollusks, worms, and small vertebrates, such as amphibians and reptiles. Larger species will eat small mammals. Feeding methods Food is taken mainly off the ground but also off the vine. The birds may jump for fruit or, as with the Crypturellus species, jump up to a metre in height for insects. The main foraging technique is a slow walk with head down, pecking at the ground and looking up occasionally. Small animals are eaten whole, larger ones are beaten against the ground or pecked. Bills rather than feet are used to probe leaf litter and sift through soil deep. The most frequent diggers are Rhynchotus, Nothura and Nothoprocta species, which are open country birds. They have their nostrils positioned at the base of the bills, a feature thought to be an adaptation to their digging. As with most birds, they swallow grit to aid their gizzards in digestion. Some species follow army ants, eating from the disturbance created. Others feed in the company of antbirds, formicariids, and ovenbirds. Nothura species, in particular, will follow livestock and eat the ticks that fall off them as well as the insects knocked off bushes as they pass. Drinking Water is required by most tinamou species, with some needing a good source within their home territory. Solitary tinamous can withstand an extended period without water by eating more succulent plants. However, species that live in arid or semi-arid climates rarely need any water additional to that ingested with their diet. When tinamous drink, unlike most other birds, they do so by sucking and swallowing, instead of lifting their heads and letting gravity do the work. Health and mortality Hygiene Tinamous are avid bathers. During heavy rain they may stand erect with their bill pointing skyward allowing the rain to wash over them. They will dust-bathe at regular intervals, and have been known to dust-bathe often enough to tint themselves the same color as the soil. They also sunbathe, and will do so while resting on one leg with an outstretched wing. Defecation for a tinamou is a slightly involved task as it must move aside the dense plumage that surrounds the cloaca to avoid soiling itself. Captive tinamous defecate once daily. Parasites There are over 240 species of bird lice that infest tinamous, with one individual bird recorded as hosting nine species. Blood parasites include louse flies, leeches, nematodes, cestodes, armadillo ticks, mites, and trematodes. Darwin's nothura may carry a malarial plasmodium. Predators Tinamou predators include cats, foxes, raccoons, skunks, weasels, tayras, rats, peccaries, and opossums. Legend speaks of jaguars that imitate the call to trick and catch them. Nests are vulnerable to snakes, monkeys and opossums. Giant anteaters have been seen on Marajo Island breaking tinamou eggs. Forest falcons and orange-breasted falcons have been seen hunting them, and vampire bats lap their blood. Relationship with humans Mythology Tinamous have established themselves in the folklore and histories of the indigenous people of South America and Central America. Forest tribes of Brazil and Colombia believe the jaguar imitates the call of the great tinamou in order to track and eat it. A tale from the Guahibo Indians tells of a young man traveling by canoe who tried to locate a calling tinamou. As he approached the bank he became suspicious at the harshness of the call and backed away just as a jaguar burst out of the vegetation. Panamanian tradition states that after the "Great Flood", the great tinamou grew frightened of the bright colors in the rainbow. He flew away from the rainbow, the ark, and the rest of the animals, heading for the darkest part of the forest, where he has remained ever since. A Brazilian legend explains of the separation of the red-winged tinamou and the undulated tinamou. The story starts off with how inseparable the two birds were, as they did everything together. One day they got into an argument and split up. The undulated tinamou went into the deepest dark of the forest, and the red-winged tinamou wanting to be different went to the grassy plains. One day, the undulated tinamou was feeling sad and lonely, went to the forest's edge and called his old friend. "Shall we make up?" he cries. The red-winged tinamou responds with "What me, never again". This story is meant to show that they are often heard but seldom seen. Introduction and translocation During the 20th century there were numerous attempts to introduce or reintroduce tinamous to various parts of the world. The red-winged tinamou has been reintroduced to the state of Rio de Janeiro, Brazil, where its wild population was hunted to extermination at the turn of the 20th century. France, Germany, and Hungary have all, unsuccessfully, attempted to introduce them into their countryside. There have been several unsuccessful attempts to introduce tinamous to the United States. In Oregon, between 1966 and 1974, 473 ornate tinamous and 110 red-winged tinamou were brought in. In 1966 and 1971, Florida introduced 128 spotted nothura. In 1969, 47 and 136 spotted nothura were introduced to Alabama and Texas. 1970 saw Colorado and Oklahoma introducing 164 and 100 Darwin's nothura respectively. In 1971, Nebraska brought 256 elegant crested tinamou, and California introduced 217 in 1969, and 1200 between 1971 and 1977. The 1885 introduction of Chilean tinamou to Easter Island was successful, though the population has not prospered since Chimango caracaras were introduced in 1928. Domestication and aviculture No tinamou species has been successfully domesticated so far, despite their ability to breed well in captivity. The red-winged tinamou has been bred on farms in France, Great Britain, Belgium, and Denmark. They, along with some of the Crypturellus species, are being bred in Rio Grande do Sul to boost numbers for hunting. Hybridization can occur. Many South American zoos hold tinamous, as do some private estates. Examples of captive breeding are small-billed tinamou in Minas Gerais and red-winged tinamou in Rio Grande do Sul. The small-billed tinamou has looked promising for domestication as the birds can raise three to four broods per year and are resistant to diseases that affect domestic chickens. Pests Several species have adapted to agricultural systems and will enter grain fields after the harvest to glean the ground; they will also enter the fields during the growing season, to the dismay of the farmers. Some, in particular the ornate and Andean tinamous, will dig up tubers such as potatoes, while red-winged tinamous create similar problems in peanut plantations. However, some members of Nothoprocta will prey on insect pests without damaging the crops. Spotted nothuras have been documented eating weeds and, of the 28 animals they were recorded as eating, 26 were considered to be pests. Research The Tinamiformes are one of the least studied orders of birds despite tinamous exhibiting rare and little understood behavioral patterns. They have male parental care which is not always associated with polyandry or sex-role reversal. Their varied mating systems and diverse habitats have the potential, through comparative studies, to explain how ecological differences affect mating strategies. In some species, females cooperate in assembling clutches of eggs for different males. However, methodological difficulties have hampered behavioral research, especially on the forest dwelling species, because their secretive behavior and cryptic coloration make them difficult to follow for continuous observation. It was not until 2003 that the first scientific symposium on tinamous was organized at the VIIth Neotropical Ornithological Congress, held at Termas de Puyehue, Chile. Status and conservation The status of the family is not easy to determine as many species live in the Amazon Basin or the far reaches of the Andes and attract little attention, even from ornithologists. Moreover, their cryptic coloration and behavior means that their presence often remains unnoticed. A large proportion of the species are Amazonian, with the majority of these decreasing in range. Most, however, are surviving well enough so far to avoid being classified as threatened. Major threats are habitat fragmentation and destruction. Although they are hunted throughout their range, it generally has little or only localized impact on the populations of the more widespread and common species. Pesticides are a problem throughout the grasslands and farmlands. The International Union for Conservation of Nature classifies seven tinamou species as vulnerable and seven as near threatened. The solitary tinamou is listed under Appendix I of CITES (the Convention on International Trade in Endangered Species of Wild Fauna and Flora). Land clearance The major threat for the forest tinamous is deforestation. Neotropical forests are badly affected, with large tracts being clearcut for cropping, pasture or timber plantations. Much of this land is poor in nutrients, so is abandoned after a few years for newly cleared land. Forest species are consequently forced to adapt, relocate, or die out. As well as the forests, most types of habitat in Middle and South America, apart from the high Andes and Patagonia, are under threat. There is controversy over the vegetative history with speculation that what is now high-altitude grassland in the Andes was once elfin forest. The dwarf tinamou is a resident of the open plains of eastern Brazil, though there are fewer than 10,000 birds left. It appears to have disappeared from its former range on the grasslands of Argentina and Paraguay. In the cerrado grasslands of Brazil the population of the lesser nothura has also decreased to fewer than 10,000 individuals because of agricultural and economic development. The practice of burning the fields is particularly detrimental to grassland tinamous. For example, the dwarf tinamou becomes intoxicated with the smoke and vulnerable to predators. Moreover, if the burning occurs during the nesting season, the eggs or chicks are roasted. The solitary tinamou, limited to the Atlantic forests of Brazil, Paraguay, and Argentina, is threatened by habitat destruction and hunting. In the cloud forests of northern South America, there are fewer than 10,000 black tinamous left. The tepui tinamou's range is limited to the tops of a handful of plateaus in the cloud forests of Venezuela, making it highly vulnerable to any threat. Hunting Tinamous have been popular game birds for many years in South America and Central America, so much so that some species' numbers have dropped. The steppe birds are more popular to hunt because they can be flushed into flight, rather than the forest birds that run to cover and hide. In the late 19th and early 20th century hunting was responsible for mass killing within the family, with the elegant crested tinamou and spotted nothura popular targets. In 1921, Argentina urged the control of commercial hunting of several bird species, including tinamou. Between 1890 and 1899, in Buenos Aires alone, 18 million tinamou were sold in meat markets. They were also marketed in North America as "South American partridge". One shipment alone comprised 360,000 birds. Frank Chapman of the American Museum of Natural History helped raise awareness about the rate of exploitation and its potential impact on the populations of the species. Eventually, the USA banned the importation of the birds. Hunting pressures remain, though at a reduced level. For example, 25,000–40,000 spotted nothura are killed annually through legal hunting, not counting poaching. Although some grassland species have increased in both range and numbers, they remain vulnerable to hunting with the use of dogs to flush them. Native peoples also are involved in the killing tinamou for meat, catching them in nooses or traps after imitating their calls. A family of seven in Ceará will consume 60 nothura per year. Tinamou species are among the most commonly harvested birds by subsistence hunting in the Americas. Some species are highly vulnerable to illegal hunting, or poaching. In Brazil, illegal hunts take place at night by torchlight. The use of flutes to imitate the calls during the breeding season to lure the birds into the open can result in local extermination. Overall, there is a lack of adequate controls in place to ensure sustainable hunting, as well as insufficient resources and determination to enforce existing regulations.	Biology and health sciences	Tinamiformes	null
207964	https://en.wikipedia.org/wiki/Hipparcos	Hipparcos	Hipparcos was a scientific satellite of the European Space Agency (ESA), launched in 1989 and operated until 1993. It was the first space experiment devoted to precision astrometry, the accurate measurement of the positions and distances of celestial objects on the sky. This permitted the first high-precision measurements of the intrinsic brightnesses, proper motions, and parallaxes of stars, enabling better calculations of their distance and tangential velocity. When combined with radial velocity measurements from spectroscopy, astrophysicists were able to finally measure all six quantities needed to determine the motion of stars. The resulting Hipparcos Catalogue, a high-precision catalogue of more than 118,200 stars, was published in 1997. The lower-precision Tycho Catalogue of more than a million stars was published at the same time, while the enhanced Tycho-2 Catalogue of 2.5 million stars was published in 2000. Hipparcos follow-up mission, Gaia, was launched in 2013. The word "Hipparcos" is an acronym for HIgh Precision PARallax COllecting Satellite and also a reference to the ancient Greek astronomer Hipparchus of Nicaea, who is noted for applications of trigonometry to astronomy and his discovery of the precession of the equinoxes. Background By the second half of the 20th century, the accurate measurement of star positions from the ground was running into essentially insurmountable barriers to improvements in accuracy, especially for large-angle measurements and systematic terms. Problems were dominated by the effects of the Earth's atmosphere, but were compounded by complex optical terms, thermal and gravitational instrument flexures, and the absence of all-sky visibility. A formal proposal to make these exacting observations from space was first put forward in 1967. The mission was originally proposed to the French space agency CNES, which considered it too complex and expensive for a single national programme and recommended that it be proposed in a multinational context. Its acceptance within the European Space Agency's scientific programme, in 1980, was the result of a lengthy process of study and lobbying. The underlying scientific motivation was to determine the physical properties of the stars through the measurement of their distances and space motions, and thus to place theoretical studies of stellar structure and evolution, and studies of galactic structure and kinematics, on a more secure empirical basis. Observationally, the objective was to provide the positions, parallaxes, and annual proper motions for some 100,000 stars with an unprecedented accuracy of 0.002 arcseconds, a target in practice eventually surpassed by a factor of two. The name of the space telescope, "Hipparcos", was an acronym for High Precision Parallax Collecting Satellite, and it also reflected the name of the ancient Greek astronomer Hipparchus, who is considered the founder of trigonometry and the discoverer of the precession of the equinoxes (due to the Earth wobbling on its axis). Satellite and payload The spacecraft carried a single all-reflective, eccentric Schmidt telescope, with an aperture of . A special beam-combining mirror superimposed two fields of view, 58° apart, into the common focal plane. This complex mirror consisted of two mirrors tilted in opposite directions, each occupying half of the rectangular entrance pupil, and providing an unvignetted field of view of about 1° × 1°. The telescope used a system of grids, at the focal surface, composed of 2688 alternate opaque and transparent bands, with a period of 1.208 arc-sec (8.2 micrometre). Behind this grid system, an image dissector tube (photomultiplier type detector) with a sensitive field of view of about 38-arc-sec diameter converted the modulated light into a sequence of photon counts (with a sampling frequency of 1200 Hz) from which the phase of the entire pulse train from a star could be derived. The apparent angle between two stars in the combined fields of view, modulo the grid period, was obtained from the phase difference of the two star pulse trains. Originally targeting the observation of some 100,000 stars, with an astrometric accuracy of about 0.002 arc-sec, the final Hipparcos Catalogue comprised nearly 120,000 stars with a median accuracy of slightly better than 0.001 arc-sec (1 milliarc-sec). An additional photomultiplier system viewed a beam splitter in the optical path and was used as a star mapper. Its purpose was to monitor and determine the satellite attitude, and in the process, to gather photometric and astrometric data of all stars down to about 11th magnitude. These measurements were made in two broad bands approximately corresponding to B and V in the (Johnson) UBV photometric system. The positions of these latter stars were to be determined to a precision of 0.03 arc-sec, which is a factor of 25 less than the main mission stars. Originally targeting the observation of around 400,000 stars, the resulting Tycho Catalogue comprised just over 1 million stars, with a subsequent analysis extending this to the Tycho-2 Catalogue of about 2.5 million stars. The attitude of the spacecraft about its center of gravity was controlled to scan the celestial sphere in a regular precessional motion maintaining a constant inclination between the spin axis and the direction to the Sun. The spacecraft spun around its Z-axis at the rate of 11.25 revolutions/day (168.75 arc-sec/s) at an angle of 43° to the Sun. The Z-axis rotated about the Sun-satellite line at 6.4 revolutions/year. The spacecraft consisted of two platforms and six vertical panels, all made of aluminum honeycomb. The solar array consisted of three deployable sections, generating around 300 W in total. Two S-band antennas were located on the top and bottom of the spacecraft, providing an omni-directional downlink data rate of 24 kbit/s. An attitude and orbit-control subsystem (comprising 5-newton hydrazine thrusters for course manoeuvres, 20-millinewton cold gas thrusters for attitude control, and gyroscopes for attitude determination) ensured correct dynamic attitude control and determination during the operational lifetime. Principles Some key features of the observations were as follows: through observations from space, the effects of astronomical seeing due to the atmosphere, instrumental gravitational flexure and thermal distortions could be obviated or minimised; all-sky visibility permitted a direct linking of the stars observed all over the celestial sphere; the two viewing directions of the satellite, separated by a large and suitable angle (58°), resulted in a rigid connection between quasi-instantaneous one-dimensional observations in different parts of the sky. In turn, this led to parallax determinations which are absolute (rather than relative, with respect to some unknown zero-point); the continuous ecliptic-based scanning of the satellite resulted in an optimum use of the available observing time, with a resulting catalogue providing reasonably homogeneous sky density and uniform astrometric accuracy over the entire celestial sphere; the various geometrical scan configurations for each star, at multiple epochs throughout the 3-year observation programme, resulted in a dense network of one-dimensional positions from which the barycentric coordinate direction, the parallax, and the object's proper motion, could be solved for in what was effectively a global least squares reduction of the totality of observations. The astrometric parameters as well as their standard errors and correlation coefficients were derived in the process; since the number of independent geometrical observations per object was large (typically of order 30) compared with the number of unknowns for the standard model (five astrometric unknowns per star), astrometric solutions not complying with this simple five-parameter model could be expanded to take into account the effects of double or multiple stars, or non-linear photocentric motions ascribed to unresolved astrometric binaries; a somewhat larger number of actual observations per object, of order 110, provided accurate and homogeneous photometric information for each star, from which mean magnitudes, variability amplitudes, and in many cases period and variability type classification could be undertaken. Development, launch and operations The Hipparcos satellite was financed and managed under the overall authority of the European Space Agency (ESA). The main industrial contractors were Matra Marconi Space (now EADS Astrium) and Alenia Spazio (now Thales Alenia Space). Other hardware components were supplied as follows: the beam-combining mirror from REOSC at Saint-Pierre-du-Perray, France; the spherical, folding and relay mirrors from Carl Zeiss AG in Oberkochen, Germany; the external straylight baffles from CASA in Madrid, Spain; the modulating grid from CSEM in Neuchâtel, Switzerland; the mechanism control system and the thermal control electronics from Dornier Satellite Systems in Friedrichshafen, Germany; the optical filters, the experiment structures and the attitude and orbit control system from Matra Marconi Space in Vélizy, France; the instrument switching mechanisms from Oerlikon-Contraves in Zürich, Switzerland; the image dissector tube and photomultiplier detectors assembled by the Dutch Space Research Organisation (SRON) in the Netherlands; the refocusing assembly mechanism designed by TNO-TPD in Delft, Netherlands; the electrical power subsystem from British Aerospace in Bristol, United Kingdom; the structure and reaction control system from Daimler-Benz Aerospace in Bremen, Germany; the solar arrays and thermal control system from Fokker Space System in Leiden, Netherlands; the data handling and telecommunications system from Saab Ericsson Space in Gothenburg, Sweden; and the apogee boost motor from SEP in France. Groups from the Institut d'Astrophysique in Liège, Belgium and the Laboratoire d'Astronomie Spatiale in Marseille, France, contributed optical performance, calibration and alignment test procedures; Captec in Dublin. Ireland, and Logica in London contributed to the on-board software and calibration. The Hipparcos satellite was launched (with the direct broadcast satellite TV-Sat 2 as co-passenger) on an Ariane 4 launch vehicle, flight V33, from Centre Spatial Guyanais, Kourou, French Guiana, on 8 August 1989. Launched into a geostationary transfer orbit (GTO), the Mage-2 apogee boost motor failed to fire, and the intended geostationary orbit was never achieved. However, with the addition of further ground stations, in addition to ESA operations control centre at European Space Operations Centre (ESOC) in Germany, the satellite was successfully operated in its geostationary transfer orbit (GTO) for almost 3.5 years. All of the original mission goals were, eventually, exceeded. Including an estimate for the scientific activities related to the satellite observations and data processing, the Hipparcos mission cost about €600 million (in year 2000 economic conditions), and its execution involved some 200 European scientists and more than 2,000 individuals in European industry. Hipparcos Input Catalogue The satellite observations relied on a pre-defined list of target stars. Stars were observed as the satellite rotated, by a sensitive region of the image dissector tube detector. This pre-defined star list formed the Hipparcos Input Catalogue (HIC): each star in the final Hipparcos Catalogue was contained in the Input Catalogue. The Input Catalogue was compiled by the INCA Consortium over the period 1982–1989, finalised pre-launch, and published both digitally and in printed form. Although fully superseded by the satellite results, it nevertheless includes supplemental information on multiple system components as well as compilations of radial velocities and spectral types which, not observed by the satellite, were not included in the published Hipparcos Catalogue. Constraints on total observing time, and on the uniformity of stars across the celestial sphere for satellite operations and data analysis, led to an Input Catalogue of some 118,000 stars. It merged two components: first, a survey of around 58,000 objects as complete as possible to the following limiting magnitudes: V<7.9 + 1.1sin\|b\| for spectral types earlier than G5, and V<7.3 + 1.1sin\|b\| for spectral types later than G5 (b is the Galactic latitude). Stars constituting this survey are flagged in the Hipparcos Catalogue. The second component comprised additional stars selected according to their scientific interest, with none fainter than about magnitude V=13 mag. These were selected from around 200 scientific proposals submitted on the basis of an Invitation for Proposals issued by ESA in 1982, and prioritised by the Scientific Proposal Selection Committee in consultation with the Input Catalogue Consortium. This selection had to balance 'a priori' scientific interest, and the observing programme's limiting magnitude, total observing time, and sky uniformity constraints. Data reductions For the main mission results, the data analysis was carried out by two independent scientific teams, NDAC and FAST, together comprising some 100 astronomers and scientists, mostly from European (ESA-member state) institutes. The analyses, proceeding from nearly 1000 Gbit of satellite data acquired over 3.5 years, incorporated a comprehensive system of cross-checking and validation, and is described in detail in the published catalogue. A detailed optical calibration model was included to map the transformation from sky to instrumental coordinates. Its adequacy could be verified by the detailed measurement residuals. The Earth's orbit, and the satellite's orbit with respect to the Earth, were essential for describing the location of the observer at each epoch of observation, and were supplied by an appropriate Earth ephemeris combined with accurate satellite ranging. Corrections due to special relativity (stellar aberration) made use of the corresponding satellite velocity. Modifications due to general relativistic light bending were significant (4 milliarc-sec at 90° to the ecliptic) and corrected for deterministically assuming γ=1 in the PPN formalism. Residuals were examined to establish limits on any deviations from this general relativistic value, and no significant discrepancies were found. Reference frame The satellite observations essentially yielded highly accurate relative positions of stars with respect to each other, throughout the measurement period (1989–1993). In the absence of direct observations of extragalactic sources (apart from marginal observations of quasar 3C 273) the resulting rigid reference frame was transformed to an inertial frame of reference linked to extragalactic sources. This allows surveys at different wavelengths to be directly correlated with the Hipparcos stars, and ensures that the catalogue proper motions are, as far as possible, kinematically non-rotating. The determination of the relevant three solid-body rotation angles, and the three time-dependent rotation rates, was conducted and completed in advance of the catalogue publication. This resulted in an accurate but indirect link to an inertial, extragalactic, reference frame. A variety of methods to establish this reference frame link before catalogue publication were included and appropriately weighted: interferometric observations of radio stars by VLBI networks, MERLIN and Very Large Array (VLA); observations of quasars relative to Hipparcos stars using charge-coupled device (CCD), photographic plates, and the Hubble Space Telescope; photographic programmes to determine stellar proper motions with respect to extragalactic objects (Bonn, Kiev, Lick, Potsdam, Yale/San Juan); and comparison of Earth rotation parameters obtained by Very-long-baseline interferometry (VLBI) and by ground-based optical observations of Hipparcos stars. Although very different in terms of instruments, observational methods and objects involved, the various techniques generally agreed to within 10 milliarc-sec in the orientation and 1 milliarc-sec/year in the rotation of the system. From appropriate weighting, the coordinate axes defined by the published catalogue are believed to be aligned with the extragalactic radio frame to within ±0.6 milliarc-sec at the epoch J1991.25, and non-rotating with respect to distant extragalactic objects to within ±0.25 milliarc-sec/yr. The Hipparcos and Tycho Catalogues were then constructed such that the resulting Hipparcos celestial reference frame (HCRF) coincides, to within observational uncertainties, with the International Celestial Reference Frame (ICRF), and representing the best estimates at the time of the catalogue completion (in 1996). The HCRF is thus a materialisation of the International Celestial Reference System (ICRS) in the optical domain. It extends and improves the J2000 (FK5) system, retaining approximately the global orientation of that system but without its regional errors. Double and multiple stars Whilst of enormous astronomical importance, double stars and multiple stars provided considerable complications to the observations (due to the finite size and profile of the detector's sensitive field of view) and to the data analysis. The data processing classified the astrometric solutions as follows: single-star solutions: 100,038 entries, of which 6,763 were flagged as suspected double component solutions (Annex C): 13,211 entries, comprising 24,588 components in 12,195 solutions acceleration solutions (Annex G): 2,622 solutions orbital solutions (Annex O): 235 entries variability-induced movers (Annex V): 288 entries stochastic solutions (Annex X): 1,561 entries no valid astrometric solution: 263 entries (of which 218 were flagged as suspected double) If a binary star has a long orbital period such that non-linear motions of the photocentre were insignificant over the short (3-year) measurement duration, the binary nature of the star would pass unrecognised by Hipparcos, but could show as a Hipparcos proper motion discrepant compared to those established from long temporal baseline proper motion programmes on ground. Higher-order photocentric motions could be represented by a 7-parameter, or even 9-parameter model fit (compared to the standard 5-parameter model), and typically such models could be enhanced in complexity until suitable fits were obtained. A complete orbit, requiring 7 elements, was determined for 45 systems. Orbital periods close to one year can become degenerate with the parallax, resulting in unreliable solutions for both. Triple or higher-order systems provided further challenges to the data processing. Photometric observations The highest accuracy photometric data were provided as a by-product of the main mission astrometric observations. They were made in a broad-band visible light passband, specific to Hipparcos, and designated Hp. The median photometric precision, for Hp<9 magnitude, was 0.0015 magnitude, with typically 110 distinct observations per star throughout the 3.5-year observation period. As part of the data reduction and catalogue production, new variables were identified and designated with appropriate variable star designations. Variable stars were classified as periodic or unsolved variables; the former were published with estimates of their period, variability amplitude, and variability type. In total some 11,597 variable objects were detected, of which 8,237 were newly classified as variable. There are, for example, 273 Cepheid variables, 186 RR Lyr variables, 108 Delta Scuti variables, and 917 eclipsing binary stars. The star mapper observations, constituting the Tycho (and Tycho-2) Catalogue, provided two colours, roughly B and V in the Johnson UBV photometric system, important for spectral classification and effective temperature determination. Radial velocities Classical astrometry concerns only motions in the plane of the sky and ignores the star's radial velocity, i.e. its space motion along the line-of-sight. Whilst critical for an understanding of stellar kinematics, and hence population dynamics, its effect is generally imperceptible to astrometric measurements (in the plane of the sky), and therefore it is generally ignored in large-scale astrometric surveys. In practice, it can be measured as a Doppler shift of the spectral lines. More strictly, however, the radial velocity does enter a rigorous astrometric formulation. Specifically, a space velocity along the line-of-sight means that the transformation from tangential linear velocity to (angular) proper motion is a function of time. The resulting effect of secular or perspective acceleration is the interpretation of a transverse acceleration actually arising from a purely linear space velocity with a significant radial component, with the positional effect proportional to the product of the parallax, the proper motion, and the radial velocity. At the accuracy levels of Hipparcos it is of (marginal) importance only for the nearest stars with the largest radial velocities and proper motions, but was accounted for in the 21 cases for which the accumulated positional effect over two years exceeds 0.1 milliarc-sec. Radial velocities for Hipparcos Catalogue stars, to the extent that they are presently known from independent ground-based surveys, can be found from the astronomical database of the Centre de données astronomiques de Strasbourg. The absence of reliable distances for the majority of stars means that the angular measurements made, astrometrically, in the plane of the sky, cannot generally be converted into true space velocities in the plane of the sky. For this reason, astrometry characterises the transverse motions of stars in angular measure (e.g. arcsec per year) rather than in km/s or equivalent. Similarly, the typical absence of reliable radial velocities means that the transverse space motion (when known) is, in any case, only a component of the complete, three-dimensional, space velocity. Published catalogues The final Hipparcos Catalogue was the result of the critical comparison and merging of the two (NDAC and FAST consortia) analyses, and contains 118,218 entries (stars or multiple stars), corresponding to an average of some three stars per square degree over the entire sky. Median precision of the five astrometric parameters (Hp<9 magnitude) exceeded the original mission goals, and are between 0.6 and 1.0 mas. Some 20,000 distances were determined to better than 10%, and 50,000 to better than 20%. The inferred ratio of external to standard errors is ≈1.0–1.2, and estimated systematic errors are below 0.1 mas. The number of solved or suspected double or multiple stars is 23,882. Photometric observations yielded multi-epoch photometry with a mean number of 110 observations per star, and a median photometric precision (Hp<9 magnitude) of 0.0015 magnitude, with 11,597 entries were identified as variable or possibly-variable. For the star mapper results, the data analysis was carried out by the Tycho Data Analysis Consortium (TDAC). The Tycho Catalogue comprises more than one million stars with 20–30 milliarc-sec astrometry and two-colour (B and V band) photometry. The final Hipparcos and Tycho Catalogues were completed in August 1996. The catalogues were published by European Space Agency (ESA) on behalf of the scientific teams in June 1997. A more extensive analysis of the star mapper (Tycho) data extracted additional faint stars from the data stream. Combined with old photographic plate observations made several decades earlier as part of the Astrographic Catalogue programme, the Tycho-2 Catalogue of more than 2.5 million stars (and fully superseding the original Tycho Catalogue) was published in 2000. The Hipparcos and Tycho-1 Catalogues were used to create the Millennium Star Atlas: an all-sky atlas of one million stars to visual magnitude 11. Some 10,000 nonstellar objects are also included to complement the catalogue data. Between 1997 and 2007, investigations into subtle effects in the satellite attitude and instrument calibration continued. A number of effects in the data that had not been fully accounted for were studied, such as scan-phase discontinuities and micrometeoroid-induced attitude jumps. A re-reduction of the associated steps of the analysis was eventually undertaken. This has led to improved astrometric accuracies for stars brighter than Hp=9.0 magnitude, reaching a factor of about three for the brightest stars (Hp<4.5 magnitude), while also underlining the conclusion that the Hipparcos Catalogue as originally published is generally reliable within the quoted accuracies. All catalogue data are available online from the Centre de données astronomiques de Strasbourg. Scientific results The Hipparcos results have affected a very broad range of astronomical research, which can be classified into three major themes: the provision of an accurate reference frame: this has allowed the consistent and rigorous re-reduction of historical astrometric measurements, including those from Schmidt plates, meridian circles, the 100-year-old Astrographic Catalogue, and 150 years of Earth-orientation measurements. These, in turn, have yielded a dense reference framework with high-accuracy, long-term proper motions (the Tycho-2 Catalogue). Reduction of current state-of-the-art survey data has yielded the dense UCAC2 Catalogue of the U.S. Naval Observatory on the same reference system, and improved astrometric data from recent surveys such as the Sloan Digital Sky Survey and 2MASS. Implicit in the high-accuracy reference frame is the measurement of gravitational lensing and the detection and characterisation of double and multiple stars; constraints on stellar structure and stellar evolution: the accurate distances and luminosities of 100,000 stars has provided the most comprehensive and accurate data set of fundamental stellar parameters to date, placing constraints on internal rotation, element diffusion, convective motions, and asteroseismology. Combined with theoretical models and other data it yields evolutionary masses, radii, and ages for large numbers of stars covering a wide range of evolutionary states; Galactic kinematics and dynamics: the uniform and accurate distances and proper motions have provided a substantial advance in understanding of stellar kinematics and the dynamical structure of the solar neighbourhood, ranging from the presence and evolution of clusters, associations and moving groups, the presence of resonance motions due to the Galaxy's central bar and spiral arms, determination of the parameters describing galactic rotation, discrimination of the disk and halo populations, evidence for halo accretion, and the measurement of space motions of runaway stars, globular clusters, and many other types of star. Associated with these major themes, Hipparcos has provided results in topics as diverse as Solar System science, including mass determinations of asteroids, Earth's rotation and Chandler wobble; the internal structure of white dwarfs; the masses of brown dwarfs; the characterisation of extra-solar planets and their host stars; the height of the Sun above the Galactic mid-plane; the age of the Universe; the stellar initial mass function and star formation rates; and strategies for the search for extraterrestrial intelligence. The high-precision multi-epoch photometry has been used to measure variability and stellar pulsations in many classes of objects. The Hipparcos and Tycho catalogues are now routinely used to point ground-based telescopes, navigate space missions, and drive public planetaria. Since 1997, several thousand scientific papers have been published making use of the Hipparcos and Tycho catalogues. A detailed review of the Hipparcos scientific literature between 1997 and 2007 was published in 2009, and a popular account of the project in 2010. Some examples of notable results include (listed chronologically): studies of Galactic rotation from Cepheid variables the nature of Delta Scuti variables studies of local stellar kinematics testing the white dwarf mass–radius relation the structure and dynamics of the Hyades cluster kinematics of Wolf–Rayet stars and O-type runaway stars subdwarf parallaxes: metal-rich clusters and the thick disk fine structure of the red giant clump and associated distance determinations unexpected stellar velocity distribution in the warped Galactic disk confirming the Lutz–Kelker bias of parallax measurement refining the Oort and Galactic constants Galactic disk dark matter, terrestrial impact cratering and the law of large numbers vertical motion and expansion of the Gould Belt the use of gamma-ray bursts as direction and time markers in SETI strategies evidence of a galaxy merger in the early formation history of the Milky Way study of nearby OB associations close approaches of stars to the Solar System studies of binary star orbits and masses the HD 209458 planetary transits formation of the stellar Galactic halo and thick disk the local density of matter in the Galaxy and the Oort limit ice age epochs and the Sun's path through the Galaxy local kinematics of K and M giants and the concept of superclusters an improved reference frame for long-term Earth rotation studies the local stellar velocity field in the Galaxy Identification of two possible "siblings" of the Sun (HIP 87382 and HIP 47399), to be studied for evidence of exoplanets The Pleiades distance controversy One controversial result has been the derived proximity, at about 120 parsecs, of the Pleiades cluster, established both from the original catalogue as well as from the revised analysis. This has been contested by various other recent work, placing the mean cluster distance at around 130 parsecs. According to a 2012 paper, the anomaly was due to the use of a weighted mean when there is a correlation between distances and distance errors for stars in clusters. It is resolved by using an unweighted mean. There is no systematic bias in the Hipparcos data when it comes to star clusters. In August 2014, the discrepancy between the cluster distance of as measured by Hipparcos and the distance of derived with other techniques was confirmed by parallax measurements made using VLBI, which gave , the most accurate and precise distance yet presented for the cluster. Polaris Another distance debate set-off by Hipparcos is for the distance to the star Polaris. Hipparcos-Gaia Hipparcos data is recently being used together with Gaia data. Especially the comparison of the proper motion of stars from both spacecraft is being used to search for hidden binary companions. Hipparcos-Gaia data is also used to measure the dynamical mass of known binaries, such as substellar companions. Hipparcos-Gaia data was used to measure the mass of the exoplanet Beta Pictoris b and is sometimes used to study other long-period exoplanets, such as HR 5183 b. People Pierre Lacroute (Observatory of Strasbourg): proposer of space astrometry in 1967 Michael Perryman: ESA project scientist (1981–1997), and project manager during satellite operations (1989–1993) Catherine Turon (Observatoire de Paris-Meudon): leader of Input Catalogue Consortium Erik Høg: leader of the TDAC Consortium Lennart Lindegren (Lund Observatory): leader of the NDAC Consortium Jean Kovalevsky: leader of the FAST Consortium Adriaan Blaauw: chair of the observing programme selection committee Hipparcos Science Team: Uli Bastian, Pierluigi Bernacca, Michel Crézé, Francesco Donati, Michel Grenon, Michael Grewing, Erik Høg, Jean Kovalevsky, Floor van Leeuwen, Lennart Lindegren, Hans van der Marel, Francois Mignard, Andrew Murray, Michael Perryman (chair), Rudolf Le Poole, Hans Schrijver, Catherine Turon Franco Emiliani: ESA project manager (1981–1985) Hamid Hassan: ESA project manager (1985–1989) Dietmar Heger: ESA/ESOC spacecraft operations manager Michel Bouffard: Matra Marconi Space project manager Bruno Strim: Alenia Spazio project manager	Technology	Space-based observatories	null
208053	https://en.wikipedia.org/wiki/Meerkat	Meerkat	The meerkat (Suricata suricatta) or suricate is a small mongoose found in southern Africa. It is characterised by a broad head, large eyes, a pointed snout, long legs, a thin tapering tail, and a brindled coat pattern. The head-and-body length is around , and the weight is typically between . The coat is light grey to yellowish-brown with alternate, poorly-defined light and dark bands on the back. Meerkats have foreclaws adapted for digging and have the ability to thermoregulate to survive in their harsh, dry habitat. Three subspecies are recognised. Meerkats are highly social, and form packs of two to 30 individuals each that occupy home ranges around in area. There is a social hierarchy—generally dominant individuals in a pack breed and produce offspring, and the nonbreeding, subordinate members provide altruistic care to the pups. Breeding occurs around the year, with peaks during heavy rainfall; after a gestation of 60 to 70 days, a litter of three to seven pups is born. They live in rock crevices in stony, often calcareous areas, and in large burrow systems in plains. The burrow systems, typically in diameter with around 15 openings, are large underground networks consisting of two to three levels of tunnels. These tunnels are around high at the top and wider below, and extend up to into the ground. Burrows have moderated internal temperatures and provide a comfortable microclimate that protects meerkats in harsh weather and at extreme temperatures. Meerkats are active during the day, mostly in the early morning and late afternoon; they remain continually alert and retreat to burrows when sensing danger. They use a broad variety of calls to communicate among one another for different purposes, for example to raise an alarm on sighting a predator. Primarily insectivorous, meerkats feed heavily on beetles and lepidopterans, arthropods, amphibians, small birds, reptiles, and plant material in their diet. Commonly living in arid, open habitats with little woody vegetation, meerkats occur in southwestern Botswana, western and southern Namibia, and northern and western South Africa; the range barely extends into southwestern Angola. With no significant threats to the population, the meerkat is listed as Least Concern on the IUCN Red List. Meerkats are widely depicted in television, movies and other media. Etymology The word 'meerkat' derives from the Dutch name for a kind of monkey, which in turn comes from the Old High German mericazza, possibly as a combination of meer ('lake') and kat ('cat'). This may be related to the similar (markat, or monkey), deriving from Sanskrit, though the Germanic origin of the word predates any known connections to India. The name was used for small mammals in South Africa from 1801 onward, possibly because the Dutch colonialists used the name in reference to many burrowing animals. The name for the meerkat is 'suricate', possibly deriving from the French 'surikate', which in turn may have a Dutch origin. In Afrikaans the meerkat is called graatjiemeerkat or stokstertmeerkat; the term mierkatte or meerkatte can refer to both the meerkat and the yellow mongoose (). In colloquial Afrikaans mier means 'ant' and kat means 'cat', hence the name probably refers to the meerkat's association with termite mounds. Taxonomy In 1776, Johann Christian Daniel von Schreber described a meerkat from the Cape of Good Hope, giving it the scientific name Viverra suricatta. The generic name Suricata was proposed by Anselme Gaëtan Desmarest in 1804, who also described a zoological specimen from the Cape of Good Hope. The present scientific name Suricata suricatta was first used by Oldfield Thomas and Harold Schwann in 1905 when they described a specimen collected at Wakkerstroom. They suggested there were four local meerkat races in the Cape and Deelfontein, Grahamstown, Orange River Colony and southern Transvaal, and Klipfontein respectively. Several zoological specimens were described between the late 18th and 20th centuries, of which three are recognised as valid subspecies: S. s. suricatta occurs in southern Namibia, southern Botswana, and South Africa. S. s. majoriae occurs in central and northwestern Namibia. S. s. iona occurs in southwestern Angola. Phylogeny and evolution Meerkat fossils dating back to have been excavated in various locations in South Africa. A 2009 phylogenetic study of the family Herpestidae suggests it split into two lineages around the Early Miocene (25.4–18.2 mya)—eusocial and solitary mongooses. The meerkat belongs to the monophyletic eusocial mongoose clade along with several other African mongooses: Crossarchus (kusimanse), Helogale (dwarf mongoose), Liberiictis (Liberian mongoose) and Mungos (banded mongoose). The solitary mongoose lineage comprises two clades including species such as Meller's mongoose (Rhynchogale melleri) and the yellow mongoose (Cynictis penicillata). The meerkat genetically diverged from the rest of the clade 22.6–15.6 mya. The phylogenetic relationships of the meerkat are depicted as follows: Characteristics The meerkat is a small mongoose of slim build characterised by a broad head, large eyes, a pointed snout, long legs, a thin tapering tail and a brindled coat pattern. It is smaller than most other mongooses except the dwarf mongooses (genus Helogale) and possibly Galerella species. The head-and-body length is around , and the weight has been recorded to be between without much variation between the sexes (though some dominant females can be heavier than the rest). The soft coat is light grey to yellowish brown with alternate, poorly-defined light and dark bands on the back. Individuals from the southern part of the range tend to be darker. The guard hairs, light at the base, have two dark rings and are tipped with black or silvery white; several such hairs aligned together give rise to the coat pattern. These hairs are typically between , but measure on the flanks. Its head is mostly white and the underparts are covered sparsely with dark reddish-brown fur, with the dark skin underneath showing through. The eyes, in sockets covering over 20% of the skull length, are capable of binocular vision. The slim, yellowish tail, unlike the bushy tails of many other mongooses, measures , and is tipped with black. Females have six nipples. The meerkat looks similar to two sympatric species—the banded and the yellow mongooses. The meerkat can be told apart from the banded mongoose by its smaller size, shorter tail and bigger eyes relative to the head; the yellow mongoose differs in having a bushy tail and lighter coat with an inner layer of yellow fur under the normal brown fur. The meerkat has 36 teeth with the dental formula of . It is well adapted for digging, movement through tunnels and standing erect, though it is not as capable of running and climbing. The big, sharp and curved foreclaws (slightly longer than the hindclaws) are highly specialised among the feliforms, and enable the meerkat to dig efficiently. The black, crescent-like ears can be closed to prevent the entry of dirt and debris while digging. The tail is used to balance when standing upright. Digitigrade, the meerkat has four digits on each foot with thick pads underneath. The meerkat has a specialised thermoregulation system that helps it survive in its harsh desert habitat. A study showed that its body temperature follows a diurnal rhythm, averaging during the day and at night. As the body temperature falls below the thermoneutral zone, determined to be , the heart rate and oxygen consumption plummet; perspiration increases sharply at temperatures above this range. Additionally, it has a basal metabolic rate remarkably lower than other carnivores, which helps in conserving water, surviving on lower amounts of food and decreasing heat output from metabolic processes. During winter, it balances heat loss by increasing the metabolic heat generation and other methods such as sunbathing. Ecology and behaviour The meerkat is a social mammal, forming packs of two to 30 individuals each comprising nearly equal numbers of either sex and multiple family units of pairs and their offspring. Members of a pack take turns at jobs such as looking after pups and keeping a lookout for predators. Meerkats are a cooperatively breeding species—typically the dominant 'breeders' in a pack produce offspring, and the nonbreeding, subordinate 'helpers' provide altruistic care for the pups. This division of labour is not as strictly defined as it is in specialised eusocial species, such as the breeder-worker distinction in ants. Moreover, meerkats have a clear dominance hierarchy with older individuals having a higher social status. A study showed that dominant individuals can contribute more to offspring care when fewer helpers were available; subordinate members increased their contributions if they could forage better. Packs live in rock crevices in stony areas and in large burrow systems in plains. A pack generally occupies a home range, large on average but sometimes as big as , containing many burrows apart, of which some remain unused. A 2019 study showed that large burrows towards the centre of a range are preferred over smaller ones located near the periphery; this was especially the case with packs that had pups to raise. A pack may shift to another burrow if the dominant female has little success finding prey in an area. The area near the periphery of home ranges is scent marked using anal gland secretions mostly by the dominant individuals; there are communal latrines, large, close to the burrows. Packs can migrate collectively in search of food, to escape high predator pressure and during floods. Meerkats are highly vigilant, and frequently survey their surroundings by turning their heads side to side; some individuals always stand sentry and look out for danger. Vocal communication is used frequently in different contexts; for instance repetitive, high-pitched barks are used to warn others of predators nearby. They will generally retreat to their burrows for safety, where they will remain until the danger is gone. They stick their heads out of burrows to check the area outside, still barking. Mobs of meerkats fiercely attack snakes that may come near them. Raptors such as bateleurs, martial eagles, tawny eagles, and pale chanting goshawks are major aerial predators; on the ground, meerkats may be threatened by bat-eared foxes, black-backed jackals, and Cape foxes. Social behaviour Encounters between members of different packs are highly aggressive, leading to severe injuries and often deaths; 19% of meerkats die by conspecific violence, which is the highest recorded percentage among mammals. Females, often the heaviest ones, try to achieve dominance over the rest in many ways such as fierce competition or taking over from the leader of the pack. A study showed that females who grew faster were more likely to assert dominance, though males did not show such a trend. Males seeking dominance over groups tend to scent mark extensively and are not submissive; they often drive out older males in a group and take over the pack themselves. Subordinate individuals face difficulties in breeding successfully; for instance, dominant females often kill the litters of subordinate ones. As such, subordinate individuals might disperse to other packs to find mates during the breeding season. Some subordinate meerkats will even kill the pups of dominant members in order to improve their own offspring's position. It can take days for emigrants to secure entry into other packs, and they often face aversion from the members. Males typically succeed in joining existing groups; they often inspect other packs and their burrow systems in search of breeding opportunities. Many often team up in 'coalitions' for as long as two months and travel nearly a day on twisted paths. Dispersal appears to be less common in females, possibly because continuing to stay within a pack can eventually win them dominance over other members. Dispersed females travel longer than coalitions, and tend to start groups of their own or join other similar females; they aim for groups of emigrant males or those without a breeding female. Subordinate females, unlike subordinate males, might be ousted from their packs, especially in the latter part of the dominant female's pregnancy, though they may be allowed to return after the birth of the pups. Burrowing Meerkat burrows are typically in diameter with around 15 openings, though one of dimensions with as many as 90 holes has been reported. These large underground networks comprise two to three levels of tunnels up to into the ground; the tunnels, around high at the top, become broader after descending around a metre. The entrances, in diameter, are created by digging at an angle of 40 degrees to the surface; the soil accumulated as a result can slightly increase the height of burrow sites. 'Boltholes' are used for a quick escape if dangers are detected. While constructing or renovating burrows meerkats will line up to form a continuous head-to-tail chain, break the soil into crumbs with their foreclaws, scoop it out with their forepaws joined and throw it behind them between their hindlegs. Outside temperatures are not reflected at once within burrows; instead there is usually an eight-hour lag which creates a temperature gradient in warrens, so that burrows are coolest in daytime and warmest at night. Temperatures inside burrows typically vary between in summer and in winter; temperatures at greater depths vary to a much lesser extent, with summer temperatures around and winter temperatures around . This reduces the need for meerkats to thermoregulate individually by providing a comfortable microclimate within burrows; moreover, burrowing protects meerkats in harsh weather and at extreme temperatures. Consequently, meerkats spend considerable time in burrows; they are active mainly during the day and return to burrows after dark and often to escape the heat of the afternoon. Activity peaks during the early morning and late afternoon. Meerkats huddle together to sleep in compact groups, sunbathe and recline on warm rocks or damp soil to adjust their body temperatures. Meerkats tend to occupy the burrows of other small mammals more than constructing them on their own; they generally share burrows with Cape ground squirrels and yellow mongooses. Cape ground squirrels and meerkats usually do not fight for space or food. Though yellow mongooses are also insectivores like meerkats, competition for prey is minimal as yellow mongooses are less selective in their diet. This association is beneficial to all the species as it saves time and efforts spent in making separate warrens. Many other species have also been recorded in the meerkat burrows, including African pygmy mice, Cape grey mongooses, four-striped grass mice, Highveld gerbils, rock hyraxes, slender mongooses, South African springhares and white-tailed rats. Vocalisations Meerkats have a broad vocal repertoire that they use to communicate among one another in several contexts; many of these calls may be combined by repetition of the same call or mixing different sounds. A study recorded 12 different types of call combinations used in different situations such as guarding against predators, caring for young, digging, sunbathing, huddling together and aggression. Short-range 'close calls' are produced while foraging and after scanning the vicinity for predators. 'Recruitment calls' can be produced to collect meerkats on sighting a snake or to investigate excrement or hair samples of predators or unfamiliar meerkats. 'Alarm calls' are given out on detecting predators. All these calls differ in their acoustic characteristics, and can evoke different responses in the 'receivers' (meerkats who hear the call); generally the greater the urgency of the scenario in which the call is given, the stronger is the response in the receivers. This indicates that meerkats are able to perceive the nature of the risk and the degree of urgency from the acoustics of a call, transmit it and respond accordingly. For instance, upon hearing a terrestrial predator alarm call, meerkats are most likely to scan the area and move towards the source of the call, while an aerial predator alarm call would most likely cause them to crouch down. A recruitment call would cause receivers to raise their tails (and often their hair) and move slowly towards the source. The complexity of calls produced by different mongooses varies by their social structure and ecology. For instance eusocial mongooses such as meerkats and banded mongooses use calls in a greater variety of contexts than do the solitary slender mongooses. Moreover, meerkats have more call types than do banded mongooses. Meerkat calls carry information to identify the signaling individual or pack, but meerkats do not appear to differentiate between calls from different sources. The calls of banded mongooses also carry a 'vocal signature' to identify the caller. Diet The meerkat is primarily an insectivore, feeding heavily on beetles and lepidopterans; it can additionally feed on eggs, amphibians, arthropods (such as scorpions, to whose venom they are immune), reptiles, small birds (such as the southern anteater-chat), plants and seeds. Captive meerkats include plenty of fruits and vegetables in their diet, and also kill small mammals by biting the backs of their skulls. They have also been observed feeding on the desert truffle Kalaharituber pfeilii. Meerkats often eat citron melons and dig out roots and tubers for their water content. Mongooses spend nearly five to eight hours foraging every day. Like other social mongooses, meerkats in a pack will disperse within of one another and browse systematically in areas within their home range without losing visual or vocal contact. Some individuals stand sentry while the rest are busy foraging. Meerkats return to an area only after a week of the last visit so that the food supply is replenished sufficiently. They hunt by scent, and often dig out soil or turn over stones to uncover hidden prey. Meerkats typically do not give chase to their prey, though they may pursue geckos and lizards over several metres. Food intake is typically low during winter. Reproduction Meerkats breed throughout the year with seasonal peaks, typically during months of heavy rainfall; for instance, maximum births occur from January to March in the southern Kalahari. Generally only dominant individuals breed, though subordinate members can also mate in highly productive years. Females become sexually mature at two to three years of age. Dominant females can have up to four litters annually (lesser for subordinate females), and the number depends on the amount of precipitation. Mating behaviour has been studied in captive individuals. Courtship behaviour is limited; the male fights with his partner, getting hold of her by her snout. He will grip the nape of her neck if she resists mounting, and hold her down by grasping her flanks during copulation. After a gestation of 60 to 70 days, a litter of three to seven pups is born. Pups weigh around in the first few days of birth; the average growth rate for the first three months is per day, typically the fastest in the first month. A 2019 study showed that growth and survival rates of pups might decrease with increase in temperature. Infants make continuous sounds that resemble bird-like tweets, that change to a shrill contact call as they grow older. Young pups are kept securely in a den, from where they emerge after around 16 days, and start foraging with adults by 26 days. The nonbreeding members of the pack help substantially with juvenile care, for instance they feed the pups and huddle with them for warmth. A study showed that nearly half of the litters of dominant females, especially those born later in the breeding season were nursed by subordinate females, mostly those that were or recently had been pregnant. Sex biases have been observed in feeding; for instance, female helpers feed female pups more than male pups unlike male helpers who feed both equally. This is possibly because the survival of female pups is more beneficial to female helpers as females are more likely to remain in their natal pack. Some helpers contribute to all activities more than others, though none of them might be specialised in any of them. Sometimes helpers favour their own needs over those of pups and decide not to feed them; this behaviour, known as "false-feeding", is more common when the prey is more valued by the meerkat. The father remains on guard and protects his offspring, while the mother spends a lot of time foraging to produce enough milk for her young. Mothers give out shrill, repetitive calls to ensure their pups follow them and remain close together. Unable to forage themselves, young pups vocalise often seeking food from their carers. Like many species, meerkat pups learn by observing and mimicking adult behaviour, though adults also engage in active instruction. For example, meerkat adults teach their pups how to eat a venomous scorpion by removing the stinger and showing the pups how to handle the creature. The mother runs around with prey in her mouth, prompting her pups to catch it. Pups become independent enough to forage at around 12 weeks of age. Meerkats are estimated to survive for five to 15 years in the wild; the maximum lifespan recorded in captivity is 20.6 years. Females appear to be able to discriminate the odour of their kin from that of others. Kin recognition is a useful ability that facilitates cooperation among relatives and the avoidance of inbreeding. When mating occurs between meerkat relatives it often results in negative fitness consequences (inbreeding depression), that affect a variety of traits such as pup mass at emergence from the natal burrow, hindleg length, growth until independence and juvenile survival. These negative effects are likely due to the increased homozygosity or higher genetic similarity among individuals that arise from inbreeding and the consequent expression of deleterious recessive mutations. Distribution and habitat The meerkat occurs in southwestern Botswana, western and southern Namibia, northern and western South Africa; the range barely extends into southwestern Angola. It lives in areas with stony, often calcareous ground in a variety of arid, open habitats with little woody vegetation. It is common in savannahs, open plains and rocky areas beside dry rivers in biomes such as the Fynbos and the Karoo, where the mean yearly rainfall is below . The average precipitation reduces to towards the northwestern areas of the range. It prefers areas with short grasses and shrubs common in velds, such as camelthorn in Namibia and Acacia in the Kalahari. It is absent from true deserts, montane regions and forests. Population densities vary greatly between places, and are significantly influenced by predators and rainfall. For instance, a study in the Kgalagadi Transfrontier Park, where predation pressure is high, recorded a lower mean meerkat density relative to a ranch with lower occurrence of predators; in response to a 10% decrease in rainfall over a year, the density fell from . Threats and conservation The meerkat is listed as Least Concern on the IUCN Red List; the population trend appears to be stable. There are no significant threats except low rainfall, which can lead to deaths of entire packs. Research has shown that temperature extremes have negative impacts on Kalahari Desert meerkats. Increased maximum air temperature is correlated with decreased survival and body mass in pups, perhaps as a result of dehydration from water loss during evaporative cooling or decreased water content in food, or from the heavier metabolic costs of thermoregulation on hot days. Higher temperatures are also associated with increased rates of endemic tuberculosis infection; this may be due to decreased immune function resulting from physiological stress, as well as increased male emigration rates observed during heat waves. Meerkats occur in several protected areas such as the Kgalagadi Transfrontier Park and the Makgadikgadi Pan. The Kalahari Meerkat Project, founded by Tim Clutton-Brock, is a long-term research project run by four different research groups that focuses on understanding cooperative behaviour in meerkats. It began in the Gemsbok National Park but was shifted to the Kuruman River Reserve in 1993. In culture Meerkats are generally tame animals. However, they are unsuitable as a pet as they can be aggressive and have a strong, ferret-like odour. In South Africa meerkats are used to kill rodents in rural households and lepidopterans in farmlands. Meerkats can transmit rabies to humans, but yellow mongooses appear to be more common vectors. It has been suggested that meerkats may even limit the spread of rabies by driving out yellow mongooses from their burrows; meerkats are generally not persecuted given their economic significance in crop protection, though they may be killed due to rabies control measures to eliminate yellow mongooses. Meerkats can also spread tick-borne diseases. Meerkats have been widely portrayed in movies, television and other media. A popular example is Timon from the Lion King franchise, who is an anthropomorphic meerkat. Meerkat Manor (2005–2008), a television programme produced by Oxford Scientific Films that was aired on Animal Planet, focused on groups of meerkats in the Kalahari that were being studied in the Kalahari Meerkat Project. Meerkats populated an acidic floating island in the 2012 film Life of Pi.	Biology and health sciences	Other carnivora	Animals
208146	https://en.wikipedia.org/wiki/Icterid	Icterid	Icterids () or New World blackbirds make up a family, the Icteridae (), of small to medium-sized, often colorful, New World passerine birds. The family contains 108 species and is divided into 30 genera. Most species have black as a predominant plumage color, often enlivened by yellow, orange, or red. The species in the family vary widely in size, shape, behavior, and coloration. The name, meaning "jaundiced ones" (from the prominent yellow feathers of many species) comes from the Ancient Greek ikteros via the Latin ictericus. This group includes the New World blackbirds, New World orioles, the bobolink, meadowlarks, grackles, cowbirds, oropendolas, and caciques. Despite the similar names, the first groups are only distantly related to the Old World common blackbird (a thrush) or the Old World orioles. The Icteridae are not to be confused with the Icteriidae, a family created in 2017 and consisting of one species — the yellow-breasted chat (Icteria virens). Characteristics Most icterid species live in the tropics, although many species also occur in temperate regions, such as the red-winged blackbird and the long-tailed meadowlark. The highest densities of breeding species are found in Colombia and southern Mexico. They inhabit a range of habitats, including scrub, swamp, forest, and savanna. Temperate species are migratory, with many species that nest in the United States and Canada moving south into Mexico and Central America. Icterids are variable in size, and often display considerable sexual dimorphism, with brighter coloration and greater size in males being typical. While such dimorphism is widely known in passerines, the sexual dimorphism by size is uniquely extreme in icterids. For example, the male great-tailed grackle is 60% heavier than the female. The smallest icterid species is the orchard oriole, in which the female averages 15 cm in length (6 in) and in weight, while the largest is the Amazonian oropendola, the male of which measures and weighs about . This variation is greater than in any other passerine family (unless the kinglet calyptura belongs with the cotingas, which would then have greater variation). One unusual morphological adaptation shared by the icterids is gaping, where the skull is configured to allow them to open their bills strongly rather than passively, allowing them to force open gaps to obtain otherwise hidden food. Most icterids have rounded tails and lack rictal bristles. They have nine primary feathers and are placed among the nine-primaried oscines. Icterids have adapted to taking a wide range of foods. Oropendolas and caciques use their gaping motion to open the skins of fruit to obtain the soft insides, and have long bills adapted to the process. Others such as cowbirds and the bobolink have shorter, stubbier bills for crushing seeds. The Jamaican blackbird uses its bill to pry amongst tree bark and epiphytes, and has adopted the evolutionary niche filled elsewhere in the Neotropics by woodcreepers. Orioles drink nectar. The nesting habits of these birds are also variable, including pendulous woven nests in the oropendolas and orioles. Many icterids are colonial, nesting in colonies of up to 100,000 birds. Some cowbird species engage in brood parasitism; females lay their eggs in the nests of other species, in a similar fashion to some cuckoos. Some species of icterid have become agricultural pests; for example, red-winged blackbirds in the United States are considered the worst vertebrate pests on some crops, such as rice. The cost of controlling blackbirds in California was $30 per acre in 1994. Not all species have been as successful, and a number of species are threatened with extinction. These include insular forms such as the Jamaican blackbird, yellow-shouldered blackbird, and St Lucia oriole, all threatened by habitat loss; and the tricolored blackbird of California, which is threatened by habitat loss and destruction of nests. Folklore Cacique and oropendola species are called paucar or similar names in Peru. As paucares are considered very intelligent, Native Americans feed the brains to their children to make them fast learners. As the male plays no part in nesting and care of the young, a man who does not work may be called a "male paucar". Taxonomy The family group was introduced in 1825 as a subfamily Icterina by Irish zoologist Nicholas Vigors. He placed the subfamily in the starling family Sturnidae. A phylogenetic analysis of the passerine families by Carl Oliveros and collaborators published in 2019 found that the family Icteridae was sister to the family Icteriidae (containing the yellow breasted chat) and together these two families formed a clade that was sister to the New World warbler family Parulidae. The genus level cladogram shown below is based on a molecular phylogenetic study by Alexis Powell and collaborators that was published in 2014. The study compared mitochondrial gene sequences. The subfamilies are those that were proposed in 2016 by Van Remsen and collaborators. The numbers of species are taken from the list maintained by Frank Gill, Pamela Rasmussen and David Donsker on behalf of the International Ornithological Committee (IOC). Genera Prehistoric icterid genera that have been described from Pleistocene fossil remains are Pandanaris from Rancho La Brea and Pyelorhamphus from Shelter Cave.	Biology and health sciences	Passerida	null
208178	https://en.wikipedia.org/wiki/New%20World%20warbler	New World warbler	The New World warblers or wood-warblers are a group of small, often colorful, passerine birds that make up the family Parulidae and are restricted to the New World. The family contains 120 species. They are not closely related to Old World warblers or Australian warblers. Most are arboreal, but some, like the ovenbird and the two waterthrushes, are primarily terrestrial. Most members of this family are insectivores. This group likely originated in northern Central America, where the greatest number of species and diversity between them is found. From there, they spread north during the interglacial periods, mainly as migrants, returning to the ancestral region in winter. Two genera, Myioborus and Basileuterus, seem to have colonized South America early, perhaps before the two continents were linked, and together constitute most warbler species of that region. The scientific name for the family, Parulidae, originates from the fact that Linnaeus in 1758 named the northern parula as a tit, Parus americanus, and as taxonomy developed, the genus name was modified first to Parulus and then to Parula. The family name derives from the name for the genus. Taxonomy The family Parulidae was introduced for the New World warblers in 1947 by American ornithologist Alexander Wetmore and collaborators with Parula as the type genus. Parula is now considered as a junior synonym of Setophaga. The family was formerly thought to be sister to a clade containing the yellow-breasted chat in its own family Icteriidae, the wrenthrush in its own family Zeledoniidae, the two Cuban warblers in the family Teretistridae and the 109 species in the family Icteridae. However, more recent studies recover them as sister to a clade containing just the yellow-breasted chat and the Icteridae, with the clade containing all three families being sister to a clade containing the chat-tanagers in Calyptophilidae, the wrenthrush, and the Phaenicophilidae. A molecular phylogenetic study of the Parulidae published in 2010 found that the species formed several major clades that did not align with the traditional genera. This led to a major reorganization of the species within the family to create monotypic genera. The changes have generally followed the recommendations of the authors of the study except in a few cases where the proposed genera were split to separate basal species from their proposed conspecifics. A large clade that included the 29 species then placed in the genus Dendroica, also included four species of Parula, one of the three species of Wilsonia and the monotypic genera Catharopeza and Setophaga. All members of the clade apart from the basal Catharopeza were placed in the expanded genus Setophaga Swainson, 1827, which under the rules of the International Code of Zoological Nomenclature, had priority over Dendroica Gray, 1842, Wilsonia Bonaparte, 1838, and Parula Bonaparte, 1838. The species that had traditionally been placed in Basileuterus formed two clades. One group retains the genus name as it includes the golden-crowned warbler, the type species for the genus. The other larger group, now with 18 species, is placed in the resurrected genus Myiothlypis Cabanis, 1850, as it contains the type species, the black-crested warbler. The genus Myioborus containing the whitestarts remained unchanged after the reorganization but six genera were no longer used: Dendroica, Ergaticus, Euthlypis, Parula, Wilsonia and Phaeothlypis. Extant Genera The family Parulidae now contains 120 species divided into 18 genera. Former species Some species that were previously placed in the Parulidae have been moved to other families: Olive warbler (Peucedramus taeniatus) – now in own family Peucedramidae Yellow-breasted chat (Icteria virens) – now in own family Icteriidae Three species in the genus Granatellus – now in the family Cardinalidae Red-breasted chat (Granatellus venustus) Grey-throated chat (Granatellus sallaei) Rose-breasted chat (Granatellus pelzelni) Wrenthrush (Zeledonia coronata) – now in own family Zeledoniidae Two species endemic to Hispaniola – now in family Phaenicophilidae Green-tailed warbler (Microligea palustris) White-winged warbler (Xenoligea montana) Two species endemic to Cuba in the genus Teretistris – now in own family Teretistridae Yellow-headed warbler (Teretistris fernandinae) Oriente warbler (Teretistris fornsi) Description All the warblers are fairly small. The smallest species is Lucy's warbler (Oreothlypis luciae), with a weight of around 6.5 g (0.23 oz) and an average length of . The Parkesia waterthrushes, the ovenbird, the russet-crowned warbler, and Semper's warbler, all of which can exceed and 21 g (0.74 oz), may be considered the largest. The migratory species tend to lay larger clutches of eggs, typically up to six, since the hazards of their journeys mean that many individuals will have only one chance to breed. In contrast, the laying of two eggs is typical for many tropical species, since the chicks can be provided with better care, and the adults are likely to have further opportunities for reproduction. Many migratory species, particularly those which breed further north, have distinctive male plumage at least in the breeding season, since males need to reclaim territory and advertise for mates each year. This tendency is particularly marked in the large genus Setophaga (formerly Dendroica). In contrast, resident tropical species, which pair for life, show little if any sexual dimorphism, but exceptions occur. The Parkesia waterthrushes and ovenbird are strongly migratory, but have identical male and female plumage, whereas the mainly tropical and sedentary yellowthroats are dimorphic. The Granatellus chats also show sexual dimorphism, but due to recent genetic work, have been moved into the family Cardinalidae (New World buntings and cardinals). The name warbler is a misnomer for the New World group of warblers established before the family was split from the Old World warbler in the 1830s. The Random House Dictionary defines "to warble" as "to sing with trills." Most New World warblers do not warble, but rather "lisp, buzz, hiss, chip, rollick, or zip."	Biology and health sciences	Passerida	null
208183	https://en.wikipedia.org/wiki/Lucky%20number	Lucky number	In number theory, a lucky number is a natural number in a set which is generated by a certain "sieve". This sieve is similar to the sieve of Eratosthenes that generates the primes, but it eliminates numbers based on their position in the remaining set, instead of their value (or position in the initial set of natural numbers). The term was introduced in 1956 in a paper by Gardiner, Lazarus, Metropolis and Ulam. In the same work they also suggested calling another sieve, "the sieve of Josephus Flavius" because of its similarity with the counting-out game in the Josephus problem. Lucky numbers share some properties with primes, such as asymptotic behaviour according to the prime number theorem; also, a version of Goldbach's conjecture has been extended to them. There are infinitely many lucky numbers. Twin lucky numbers and twin primes also appear to occur with similar frequency. However, if Ln denotes the n-th lucky number, and pn the n-th prime, then Ln > pn for all sufficiently large n. Because of their apparent similarities with the prime numbers, some mathematicians have suggested that some of their common properties may also be found in other sets of numbers generated by sieves of a certain unknown form, but there is little theoretical basis for this conjecture. The sieving process Continue removing the nth remaining numbers, where n is the next number in the list after the last surviving number. Next in this example is 9. One way that the application of the procedure differs from that of the Sieve of Eratosthenes is that for n being the number being multiplied on a specific pass, the first number eliminated on the pass is the n-th remaining number that has not yet been eliminated, as opposed to the number 2n. That is to say, the list of numbers this sieve counts through is different on each pass (for example 1, 3, 7, 9, 13, 15, 19... on the third pass), whereas in the Sieve of Eratosthenes, the sieve always counts through the entire original list (1, 2, 3...). When this procedure has been carried out completely, the remaining integers are the lucky numbers (those that happen to be prime are in bold): 1, 3, 7, 9, 13, 15, 21, 25, 31, 33, 37, 43, 49, 51, 63, 67, 69, 73, 75, 79, 87, 93, 99, 105, 111, 115, 127, 129, 133, 135, 141, 151, 159, 163, 169, 171, 189, 193, 195, 201, 205, 211, 219, 223, 231, 235, 237, 241, 259, 261, 267, 273, 283, 285, 289, 297, 303, 307, 319, 321, 327, 331, 339, ... . The lucky number which removes n from the list of lucky numbers is: (0 if n is a lucky number) 0, 2, 0, 2, 3, 2, 0, 2, 0, 2, 3, 2, 0, 2, 0, 2, 3, 2, 7, 2, 0, 2, 3, 2, 0, 2, 9, 2, 3, 2, 0, 2, 0, 2, 3, 2, 0, 2, 7, 2, 3, 2, 0, 2, 13, 2, 3, 2, 0, 2, 0, 2, 3, 2, 15, 2, 9, 2, 3, 2, 7, 2, 0, 2, 3, 2, 0, 2, 0, 2, 3, 2, 0, 2, 0, 2, 3, 2, 0, 2, 7, 2, 3, 2, 21, 2, ... Lucky primes A "lucky prime" is a lucky number that is prime. They are: 3, 7, 13, 31, 37, 43, 67, 73, 79, 127, 151, 163, 193, 211, 223, 241, 283, 307, 331, 349, 367, 409, 421, 433, 463, 487, 541, 577, 601, 613, 619, 631, 643, 673, 727, 739, 769, 787, 823, 883, 937, 991, 997, ... . It has been conjectured that there are infinitely many lucky primes.	Mathematics	Sequences	null
208186	https://en.wikipedia.org/wiki/Dione%20%28moon%29	Dione (moon)	Dione (), also designated Saturn IV, is the fourth-largest moon of Saturn. With a mean diameter of 1,123 km and a density of about 1.48 g/cm3, Dione is composed of an icy mantle and crust overlying a silicate rocky core, with rock and water ice roughly equal in mass. Its trailing hemisphere is marked by large cliffs and scarps called chasmata; the trailing hemisphere is also significantly darker compared to the leading hemisphere. The moon was discovered by Italian astronomer Giovanni Domenico Cassini in 1684 and is named after the Titaness Dione in Greek mythology. Dione was first imaged up-close by the Voyager 1 space probe in 1980. Later, the Cassini spacecraft made multiple flybys of Dione throughout the 2000s and 2010s as part of its campaign to explore the Saturn system. Name Giovanni Domenico Cassini named the four moons he discovered (Tethys, Dione, Rhea, and Iapetus) Sidera Lodoicea ("the stars of Louis") to honor king Louis XIV. Cassini found Dione in 1684 using a large aerial telescope he set up on the grounds of the Paris Observatory. The satellites of Saturn were not named until 1847, when William Herschel's son John Herschel published Results of Astronomical Observations made at the Cape of Good Hope, suggesting that the names of the Titans (sisters and brothers of Cronus) be used. Orbit Dione orbits Saturn with a semimajor axis about 2% less than that of the Moon. However, reflecting Saturn's greater mass (95 times that of Earth), Dione's orbital period is one tenth that of the Moon. Dione is currently in a 1:2 mean-motion orbital resonance with moon Enceladus, completing one orbit of Saturn for every two orbits completed by Enceladus. This resonance maintains Enceladus's orbital eccentricity (0.0047), providing a source of heat for Enceladus's extensive geological activity, which shows up most dramatically in its cryovolcanic geyser-like jets. The resonance also maintains a smaller eccentricity in Dione's orbit (0.0022), tidally heating it as well. Trojans Dione has two co-orbital, or trojan, moons, Helene and Polydeuces. They are located within Dione's Lagrangian points and , 60 degrees ahead of and behind Dione respectively. A leading co-orbital moon twelve degrees ahead of Helene was reported by Stephen P. Synnott in 1982. Physical characteristics and interior At in diameter, Dione is the 15th largest moon in the Solar System, and is more massive than all known moons smaller than itself combined. It is also Saturn's fourth-largest moon. Based on its density, Dione’s interior is likely a combination of silicate rock and water ice in nearly equal parts by mass. Shape and gravity observations collected by Cassini suggest a roughly 400 km radius rocky core surrounded by a roughly 160 km envelope of H2O, mainly in the form of water ice, but with some models suggesting that the lowermost part of this layer could be in the form of an internal liquid salt water ocean (a situation similar to that of its orbital resonance partner, Enceladus). Downward bending of the surface associated with the 1.5 km high ridge Janiculum Dorsa can most easily be explained by the presence of such an ocean. Neither moon has a shape close to hydrostatic equilibrium; the deviations are maintained by isostasy. Dione's ice shell is thought to vary in thickness by less than 5%, with the thinnest areas at the poles, where tidal heating of the crust is greatest. Though somewhat smaller and denser, Dione is otherwise very similar to Rhea. They both have similar albedo features and varied terrain, and both have dissimilar leading and trailing hemispheres. Dione's leading hemisphere is heavily cratered and is uniformly bright. Its trailing hemisphere, however, contains an unusual and distinctive surface feature: a network of bright ice cliffs. Scientists recognise Dionean geological features of the following types: Chasmata (chasms; long, deep, steep-sided depressions or canyons) Dorsa (ridges) Fossae (long narrow depressions) Craters Catenae (crater chains) Ice cliffs (formerly 'wispy terrain') When the Voyager space probe photographed Dione in 1980, it showed what appeared to be wispy features covering its trailing hemisphere. The origin of these features was mysterious, because all that was known was that the material has a high albedo and is thin enough that it does not obscure the surface features underneath. One hypothesis was that shortly after its formation Dione was geologically active, and some process such as cryovolcanism resurfaced much of its surface, with the streaks forming from eruptions along cracks in the Dionean surface that fell back as snow or ash. Later, after the internal activity and resurfacing ceased, cratering continued primarily on the leading hemisphere and wiped out the streak patterns there. This hypothesis was proven wrong by the Cassini probe flyby of 13 December 2004, which produced close-up images. These revealed that the 'wisps' were, in fact, not ice deposits at all, but rather bright ice cliffs created by tectonic fractures (chasmata). Dione has been revealed as a world riven by enormous fractures on its trailing hemisphere. The Cassini orbiter performed a closer flyby of Dione at on 11 October 2005, and captured oblique images of the cliffs, showing that some of them are several hundred metres high. Linear features Dione features linear 'virgae' that are up to hundreds of km long but less than 5 km wide. These lines run parallel to the equator and are only apparent at lower latitudes (at less than 45° north or south); similar features are noted on Rhea. They are brighter than everything around them and appear to overlay other features such as ridges and craters, indicating they are relatively young. It has been proposed that these lines are of exogenic origin, as the result of the emplacement of material across the surface by low‐velocity impacts of material sourced from Saturn's rings, co‐orbital moons, or closely approaching comets. Craters Dione's icy surface includes heavily cratered terrain, moderately cratered plains, lightly cratered plains, and areas of tectonic fractures. The heavily cratered terrain has numerous craters greater than in diameter. The plains areas tend to have craters less than in diameter. Some of the plains are more heavily cratered than others. Much of the heavily cratered terrain is located on the trailing hemisphere, with the less cratered plains areas present on the leading hemisphere. This is the opposite of what some scientists expected; Shoemaker and Wolfe proposed a cratering model for a tidally locked satellite with the highest cratering rates on the leading hemisphere and the lowest on the trailing hemisphere. This suggests that during the period of heavy bombardment, Dione was tidally locked to Saturn in the opposite orientation. Because Dione is relatively small, an impact causing a 35 kilometer crater could have spun the satellite. Because there are many craters larger than , Dione could have been repeatedly spun during its early heavy bombardment. The pattern of cratering since then and the bright albedo of the leading side suggests that Dione has remained in its current orientation for several billion years. Like Callisto, Dione's craters lack the high-relief features seen on the Moon and Mercury; this is probably due to slumping of the weak icy crust over geologic time. Atmosphere On 7 April 2010, instruments on board the uncrewed Cassini probe, which flew by Dione, detected a thin layer of molecular oxygen ions () around Dione, so thin that scientists prefer to call it an exosphere rather than a tenuous atmosphere. The density of molecular oxygen ions determined from the Cassini plasma spectrometer data ranges from 0.01 to 0.09 per cm3. The Cassini probe instruments were unable to directly detect water from the exosphere due to high background levels, but it seems that highly charged particles from the planet's powerful radiation belts could split the water in the ice into hydrogen and oxygen. Exploration Dione was first imaged by the Voyager space probes. It has also been probed five times from close distances by the Cassini orbiter. There was a close targeted flyby at a distance of on 11 October 2005; another flyby was performed on 7 April 2010, also at a distance of 500 km. A third flyby was performed on 12 December 2011 at a distance of . The following flyby was on 16 June 2015 at a distance of , and the last Cassini flyby was performed on 17 August 2015 at a distance of . In May 2013, it was announced that NASA's spacecraft Cassini had provided scientists with evidence that Dione is more active than previously realized. Using topographic data, NASA teams deduced that crustal depression associated with a prominent mountain ridge on the leading hemisphere is best explained if there was a global subsurface liquid ocean like that of Enceladus. The ridge Janiculum Dorsa has a height of 1 to 2 km (0.6 to 1.2 miles); Dione's crust seems to pucker 0.5 km (0.3 miles) under it, suggesting that the icy crust was warm when the ridge formed, probably due to the presence of a subsurface liquid ocean, which increases tidal flexing.	Physical sciences	Solar System	Astronomy
208207	https://en.wikipedia.org/wiki/Shoebill	Shoebill	The shoebill (Balaeniceps rex), also known as the whale-headed stork, and shoe-billed stork, is a large long-legged wading bird. It derives its name from its enormous shoe-shaped bill. It has a somewhat stork-like overall form and has previously been classified with the storks in the order Ciconiiformes based on this morphology. However, genetic evidence places it with pelicans and herons in the Pelecaniformes. The adult is mainly grey while the juveniles are more brown. It lives in tropical East Africa in large swamps from South Sudan to Zambia. Taxonomy The shoebill may have been known to Ancient Egyptians but was not classified until the 19th century, after skins and eventually live specimens were brought to Europe. John Gould very briefly described it in 1850 from the skin of a specimen collected on the upper White Nile by the English traveller Mansfield Parkyns. Gould provided a more detailed description in the following year. He placed the species in its own genus Balaeniceps and coined the binomial name Balaeniceps rex. The genus name comes from the Latin words balaena "whale", and caput "head", abbreviated to -ceps in compound words. Alternative common names are whalebill, shoe-billed stork and whale-headed stork. Traditionally considered as allied with the storks (Ciconiiformes), it was retained there in the Sibley-Ahlquist taxonomy which lumped a massive number of unrelated taxa into their "Ciconiiformes". Based on osteological evidence, the suggestion of a pelecaniform affinity was made in 1957 by Patricia Cottam. Microscopic analysis of eggshell structure by Konstantin Mikhailov in 1995 found that the eggshells of shoebills closely resembled those of other Pelecaniformes in having a covering of thick microglobular material over the crystalline shells. In 2003, the shoebill was again suggested as closer to the pelicans (based on anatomical comparisons) or the herons (based on biochemical evidence). A 2008 DNA study reinforces their membership among the Pelecaniformes. So far, two fossilized relatives of the shoebill have been described: Goliathia from the Early Oligocene of Egypt and Paludiavis from the Late Miocene of Pakistan and Tunisia. It has been suggested that the enigmatic African fossil bird Eremopezus have features resembling those of the shoebill and the secretary bird. Description The shoebill is a tall bird, with a typical height range of and some specimens reaching as much as . Length from tail to beak can range from and wingspan is . Weight has reportedly ranged from . A male will weigh on average around and is larger than a typical female of . The signature feature of the species is its huge, bulbous bill, which is straw-coloured with erratic greyish markings. The exposed culmen (or the measurement along the top of the upper mandible) is , the third longest bill among extant birds after pelicans and large storks, and can outrival the pelicans in bill circumference, especially if the bill is considered as the hard, bony keratin portion. As in the pelicans, the upper mandible is strongly keeled, ending in a sharp nail. The dark coloured legs are fairly long, with a tarsus length of . The shoebill's feet are exceptionally large, with the middle toe reaching in length, likely assisting the species in its ability to stand on aquatic vegetation while hunting. The neck is relatively shorter and thicker than other long-legged wading birds such as herons and cranes. The wings are broad, with a wing chord length of , and well-adapted to soaring. The plumage of adult birds is blue-grey with darker slaty-grey flight feathers. The breast presents some elongated feathers, which have dark shafts. The juvenile has a similar plumage colour, but is a darker grey with a brown tinge. When they are first born, shoebills have a more modestly-sized bill, which is initially silvery-grey. The bill becomes more noticeably large when the chicks are 23 days old and becomes well developed by 43 days. Voice The shoebill is normally silent, but they perform bill-clattering displays at the nest. When engaging in these displays, adult birds have also been noted to utter a cow-like moo as well as high-pitched whines. Both nestlings and adults engage in bill-clattering during the nesting season as a means of communication. When young are begging for food, they call out with a sound uncannily like human hiccups. In one case, a flying adult bird was heard uttering hoarse croaks, apparently as a sign of aggression at a nearby marabou stork (Leptoptilos crumeniferus). Flight pattern Its wings are held flat while soaring and, as in the pelicans and the storks of the genus Leptoptilos, the shoebill flies with its neck retracted. Its flapping rate, at an estimated 150 flaps per minute, is one of the slowest of any bird, with the exception of the larger stork species. The pattern is alternating flapping and gliding cycles of approximately seven seconds each, putting its gliding distance somewhere between the larger storks and the Andean condor (Vultur gryphus). When flushed, shoebills usually try to fly no more than . Long flights of the shoebill are rare, and only a few flights beyond its minimum foraging distance of have been recorded. Distribution and habitat The shoebill is distributed in freshwater swamps of central tropical Africa, from southern Sudan and South Sudan through parts of eastern Congo, Rwanda, Uganda, western Tanzania and northern Zambia. The species is most numerous in the West Nile sub-region and South Sudan (especially the Sudd, a main stronghold for the species); it is also significant in wetlands of Uganda and western Tanzania. More isolated records have been reported of shoebills in Kenya, the Central African Republic, northern Cameroon, south-western Ethiopia, and Malawi. Vagrant strays to the Okavango Basin, Botswana and the upper Congo River have also been sighted. The distribution of this species seems to largely coincide with that of papyrus and lungfish. They are often found in areas of flood plain interspersed with undisturbed papyrus and reedbeds. When shoebill storks are in an area with deep water, a bed of floating vegetation is a requirement. They are also found where there is poorly oxygenated water. This causes the fish living in the water to surface for air more often, increasing the likelihood a shoebill stork will successfully capture it. The shoebill is non-migratory with limited seasonal movements due to habitat changes, food availability and disturbance by humans. Petroglyphs from Oued Djerat, eastern Algeria, show that the shoebill occurred during the Early Holocene much more to the north, in the wetlands that covered the present-day Sahara Desert at that time. The shoebill occurs in extensive, dense freshwater marshes. Almost all wetlands that attract the species have undisturbed Cyperus papyrus and reed beds of Phragmites and Typha. Although their distribution largely seems to correspond with the distribution of papyrus in central Africa, the species seems to avoid pure papyrus swamps and is often attracted to areas with mixed vegetation. More rarely, the species has been seen foraging in rice fields and flooded plantations. Behaviour and ecology The shoebill is noted for its slow movements and tendency to stay still for long periods, resulting in descriptions of the species as "statue-like". They are quite sensitive to human disturbance and may abandon their nests if flushed by humans. However, while foraging, if dense vegetation stands between it and humans, this wader can be fairly tame. The shoebill is often attracted to poorly oxygenated waters such as swamps, marshes, and bogs where fish frequently surface to breathe. They also seem to exhibit migratory behaviors based upon differences in the surface water level. Immature shoebills abandon nesting sites which increased in the surface water level whereas adult shoebills abandon nesting sites which decreased in surface water level. It is suggested that both adult and immature shoebills prefer nesting sites with similar surface water levels. Exceptionally for a bird this large, the shoebill often stands and perches on floating vegetation, making them appear somewhat like a giant jacana, although the similarly sized and occasionally sympatric Goliath heron (Ardea goliath) is also known to stand on aquatic vegetation. Shoebills, being solitary, forage at or more from one another even where relatively densely populated. This species stalks its prey patiently, in a slow and lurking fashion. While hunting, the shoebill strides very slowly and is frequently motionless. Unlike some other large waders, this species hunts entirely using vision and is not known to engage in tactile hunting. When prey is spotted, it launches a quick violent strike. However, depending on the size of the prey, handling time after the strike can exceed 10 minutes. Around 60% of strikes yield prey. Frequently water and vegetation is snatched up during the strike and is spilled out from the edges of the mandibles. The activity of hippopotamus may inadvertently benefit the shoebill, as submerged hippos occasionally force fish to the surface. Breeding The solitary nature of shoebills extends to their breeding habits. Nests typically occur at less than three nests per square kilometre, unlike herons, cormorants, pelicans, and storks, which predominantly nest in colonies. The breeding pair of shoebills vigorously defends a territory of from conspecifics. In the extreme north and south of the species' range, nesting starts right after the rains end. In more central regions of the range, it may nest near the end of the wet season in order for the eggs to hatch around the beginning of the following wet season. Both parents engage in building the nest on a floating platform after clearing out an area of approximately across. The large, flattish nesting platform is often partially submerged in water and can be as much as deep. The nest itself is about wide. Both the nest and platform are made of aquatic vegetation. From one to three white eggs are laid. These eggs measure high by and weigh around . Incubation lasts for approximately 30 days. Both parents actively brood, shade, guard and feed the nestling, though the females are perhaps slightly more attentive. Shoebills use their mandibles to cool their eggs with water during days with high temperatures around . They fill their mandible once, swallow the water, and fill another mandible full of water before proceeding back to their nest where they pour out the water and regurgitate the previously swallowed water onto both the nest and eggs. Food items are regurgitated whole from the gullet straight into the bill of the young. Shoebills rarely raise more than one chick but will hatch more. The younger chicks usually die and are intended as "back-ups" in case the eldest chick dies or is weak. Fledging is reached at around 105 days and the young birds can fly well by 112 days. However, they are still fed for possibly a month or more after this. It will take the young shoebills three years before they become fully sexually mature. Shoebills are elusive when nesting, so cameras must be placed to observe them from afar to collect behavioral data. There is an advantage for birds that are early breeders, as the chicks are tended for a longer period. Diet Shoebills are largely piscivorous but are assured predators of a considerable range of wetland vertebrates. Preferred prey species have reportedly included marbled lungfish (Protopterus aethiopicus), African lungfish (Protopterus annectens), and Senegal bichir (Polypterus senegalus), various Tilapia species and catfish, the latter mainly in the genus Clarias. Other prey eaten by this species has included frogs, water snakes, Nile monitors (Varanus niloticus) and baby crocodiles. More rarely, small turtles, snails, rodents, small waterfowl and carrion have reportedly been eaten. Given its sharp-edged beak, huge bill, and wide gape, the shoebill can hunt large prey, often targeting prey bigger than is taken by other large wading birds. In the Bangweulu Swamps of Zambia, fish eaten by this species are commonly in the range of . The main prey items fed to young by the parents were the catfish Clarias gariepinus, (syn. C. mossambicus) and long water snakes. In Uganda, lungfish and catfish were mainly fed to the young. Larger lungfish and catfish were taken in Malagarasi wetlands in western Tanzania. During this study, fish around were quite frequently taken and the largest fish caught by the shoebill was 99 cm long. Fish exceeding 60 cm were usually cut into sections and swallowed at intervals. The entire process from scooping to swallowing ranged from 2 to 30 minutes depending on prey size. However, these large prey are relatively hard to handle and often targeted by African fish eagle (Icthyophaga vocifer), which frequently steal large wading bird's prey. Relationship to humans This species is considered to be one of the five most desirable birds in Africa by birdwatchers. They are docile with humans and show no threatening behavior. Researchers were able to observe a bird on its nest at a close distance – within . Shoebills are often kept in zoos, but breeding is rarely reported. Shoebills have bred successfully at Pairi Daiza in Belgium and at Tampa's Lowry Park Zoo in Florida. Appearances in popular culture Beginning in 2014 and with various interspersed surges of attention since then, the shoebill has become the subject of internet memes, in part due to its intimidating appearance and its tendency to stand still for long periods of time. One such example is a video of a shoebill standing in the rain whilst staring into the camera. These memes have since also appeared on the social media platform TikTok, bringing a comparatively unknown species of bird into popular culture. The shoebill also inspired the design of the Loftwing birds in the 2011 game The Legend of Zelda: Skyward Sword. Status and conservation The population is estimated at between and individuals, the majority of which live in swamps in South Sudan, Uganda, eastern Democratic Republic of the Congo, and Zambia. There is also a viable population in the Malagarasi wetlands in Tanzania. BirdLife International has classified it as Vulnerable with the main threats being habitat destruction, disturbance and hunting. The bird is listed under Appendix II of the Convention on International Trade in Endangered Species of Wild Fauna and Flora (CITES). Habitat destruction and degradation, hunting, disturbance and illegal capture are all contributing factors to the decline of this species. Agriculture cultivation and pasture for cattle have also caused significant habitat loss. Indigenous communities that surround Shoebill habitats capture their eggs and chicks for human consumption and for trade. Frequent fires in southern Sudan and deliberate fires for grazing access contribute to habitat loss. Some swamps in Sudan are being drained for construction of a canal to control nearby waterways, causing more habitat loss.	Biology and health sciences	Pelecanimorphae	null
208215	https://en.wikipedia.org/wiki/Geochronology	Geochronology	Geochronology is the science of determining the age of rocks, fossils, and sediments using signatures inherent in the rocks themselves. Absolute geochronology can be accomplished through radioactive isotopes, whereas relative geochronology is provided by tools such as paleomagnetism and stable isotope ratios. By combining multiple geochronological (and biostratigraphic) indicators the precision of the recovered age can be improved. Geochronology is different in application from biostratigraphy, which is the science of assigning sedimentary rocks to a known geological period via describing, cataloging and comparing fossil floral and faunal assemblages. Biostratigraphy does not directly provide an absolute age determination of a rock, but merely places it within an interval of time at which that fossil assemblage is known to have coexisted. Both disciplines work together hand in hand, however, to the point where they share the same system of naming strata (rock layers) and the time spans utilized to classify sublayers within a stratum. The science of geochronology is the prime tool used in the discipline of chronostratigraphy, which attempts to derive absolute age dates for all fossil assemblages and determine the geologic history of the Earth and extraterrestrial bodies. Dating methods Radiometric dating By measuring the amount of radioactive decay of a radioactive isotope with a known half-life, geologists can establish the absolute age of the parent material. A number of radioactive isotopes are used for this purpose, and depending on the rate of decay, are used for dating different geological periods. More slowly decaying isotopes are useful for longer periods of time, but less accurate in absolute years. With the exception of the radiocarbon method, most of these techniques are actually based on measuring an increase in the abundance of a radiogenic isotope, which is the decay-product of the radioactive parent isotope. Two or more radiometric methods can be used in concert to achieve more robust results. Most radiometric methods are suitable for geological time only, but some such as the radiocarbon method and the 40Ar/39Ar dating method can be extended into the time of early human life and into recorded history. Some of the commonly used techniques are: Radiocarbon dating. This technique measures the decay of carbon-14 in organic material and can be best applied to samples younger than about 60,000 years. Uranium–lead dating. This technique measures the ratio of two lead isotopes (lead-206 and lead-207) to the amount of uranium in a mineral or rock. Often applied to the trace mineral zircon in igneous rocks, this method is one of the two most commonly used (along with argon–argon dating) for geologic dating. Monazite geochronology is another example of U–Pb dating, employed for dating metamorphism in particular. Uranium–lead dating is applied to samples older than about 1 million years. Uranium–thorium dating. This technique is used to date speleothems, corals, carbonates, and fossil bones. Its range is from a few years to about 700,000 years. Potassium–argon dating and argon–argon dating. These techniques date metamorphic, igneous and volcanic rocks. They are also used to date volcanic ash layers within or overlying paleoanthropologic sites. The younger limit of the argon–argon method is a few thousand years. Electron spin resonance (ESR) dating Fission-track dating Cosmogenic nuclide geochronology A series of related techniques for determining the age at which a geomorphic surface was created (exposure dating), or at which formerly surficial materials were buried (burial dating). Exposure dating uses the concentration of exotic nuclides (e.g. 10Be, 26Al, 36Cl) produced by cosmic rays interacting with Earth materials as a proxy for the age at which a surface, such as an alluvial fan, was created. Burial dating uses the differential radioactive decay of 2 cosmogenic elements as a proxy for the age at which a sediment was screened by burial from further cosmic rays exposure. Luminescence dating Luminescence dating techniques observe 'light' emitted from materials such as quartz, diamond, feldspar, and calcite. Many types of luminescence techniques are utilized in geology, including optically stimulated luminescence (OSL), cathodoluminescence (CL), and thermoluminescence (TL). Thermoluminescence and optically stimulated luminescence are used in archaeology to date 'fired' objects such as pottery or cooking stones and can be used to observe sand migration. Incremental dating Incremental dating techniques allow the construction of year-by-year annual chronologies, which can be fixed (i.e. linked to the present day and thus calendar or sidereal time) or floating. Dendrochronology Ice cores Lichenometry Varves Paleomagnetic dating A sequence of paleomagnetic poles (usually called virtual geomagnetic poles), which are already well defined in age, constitutes an apparent polar wander path (APWP). Such a path is constructed for a large continental block. APWPs for different continents can be used as a reference for newly obtained poles for the rocks with unknown age. For paleomagnetic dating, it is suggested to use the APWP in order to date a pole obtained from rocks or sediments of unknown age by linking the paleopole to the nearest point on the APWP. Two methods of paleomagnetic dating have been suggested: (1) the angular method and (2) the rotation method. The first method is used for paleomagnetic dating of rocks inside of the same continental block. The second method is used for the folded areas where tectonic rotations are possible. Magnetostratigraphy Magnetostratigraphy determines age from the pattern of magnetic polarity zones in a series of bedded sedimentary and/or volcanic rocks by comparison to the magnetic polarity timescale. The polarity timescale has been previously determined by dating of seafloor magnetic anomalies, radiometrically dating volcanic rocks within magnetostratigraphic sections, and astronomically dating magnetostratigraphic sections. Chemostratigraphy Global trends in isotope compositions, particularly carbon-13 and strontium isotopes, can be used to correlate strata. Correlation of marker horizons Marker horizons are stratigraphic units of the same age and of such distinctive composition and appearance that, despite their presence in different geographic sites, there is certainty about their age-equivalence. Fossil faunal and floral assemblages, both marine and terrestrial, make for distinctive marker horizons. Tephrochronology is a method for geochemical correlation of unknown volcanic ash (tephra) to geochemically fingerprinted, dated tephra. Tephra is also often used as a dating tool in archaeology, since the dates of some eruptions are well-established. Geological hierarchy of chronological periodization Geochronology, from largest to smallest: Supereon Eon Era Period Epoch Age Chron Differences from chronostratigraphy It is important not to confuse geochronologic and chronostratigraphic units. Geochronological units are periods of time, thus it is correct to say that Tyrannosaurus rex lived during the Late Cretaceous Epoch. Chronostratigraphic units are geological material, so it is also correct to say that fossils of the genus Tyrannosaurus have been found in the Upper Cretaceous Series. In the same way, it is entirely possible to go and visit an Upper Cretaceous Series deposit – such as the Hell Creek deposit where the Tyrannosaurus fossils were found – but it is naturally impossible to visit the Late Cretaceous Epoch as that is a period of time.	Physical sciences	Geological history	null
208296	https://en.wikipedia.org/wiki/Screamer	Screamer	The screamers are three South American bird species placed in family Anhimidae. They were thought to be related to the Galliformes because of similar bills, but are more closely related to the family Anatidae, i.e. ducks and allies, and the magpie goose, within the clade Anseriformes. The clade is exceptional within the living birds in lacking uncinate processes of ribs. The three species are: The horned screamer (Anhima cornuta); the southern screamer or crested screamer (Chauna torquata); and the northern screamer or black-necked screamer (Chauna chavaria). Systematics and evolution Screamers have a poor fossil record. A putative Eocene specimen is known from Wyoming, while the more modern Chaunoides antiquus is known from the late Oligocene to early Miocene in Brazil. Anhimids are most similar to presbyornithids, with which they form a clade to the exclusion of the rest of Anseriformes. Given the presence of lamelae in the otherwise fowl-like beaks of screamers, it is even possible that they evolved from presbyornithid-grade birds, reverting from a filter-feeding lifestyle to an herbivorous one. Distribution and habitat The three species occur only in South America, ranging from Colombia to northern Argentina. The horned screamer was once present on the Caribbean island of Trinidad, but is now extirpated from there. They are large, bulky birds, with a small downy head, long legs and large feet which are only partially webbed. They have large spurs on their wings which are used in fights over mates and territorial disputes; these can break off in the breast of other screamers, and are regularly renewed. Unlike ducks, they have a partial moult and are able to fly throughout the year. They live in open areas and marshes with some grass and feed on water plants. One species, the southern screamer, is considered a pest as it raids crops and competes with farm birds. Behaviour and ecology Screamers typical lay 4–5 white eggs, with clutches ranging between 2 and 7 . Like most Anseriformes, the chicks can run as soon as they are hatched. They can swim better than they can run, so young screamers are usually raised in or near water, where they can better avoid predators. Like ducks, screamer chicks imprint early in life. That, and their unspecialized omnivorous diet makes them amenable to domestication. They can be excellent guard animals, due to their loud alarm calls ("screams") when encountering anything new and possibly threatening. Status and conservation Both the southern and the horned screamer remain widespread and are overall fairly common. In contrast, the northern screamer is relatively rare and consequently considered near threatened. They are seldom hunted, in spite of their conspicuous nature, because their flesh has a spongy texture and is riddled with air-sacs, making it highly unpalatable. The main threats are habitat destruction and increased intensification of agriculture.	Biology and health sciences	Anseriformes	Animals
208303	https://en.wikipedia.org/wiki/Primary%20production	Primary production	In ecology, primary production is the synthesis of organic compounds from atmospheric or aqueous carbon dioxide. It principally occurs through the process of photosynthesis, which uses light as its source of energy, but it also occurs through chemosynthesis, which uses the oxidation or reduction of inorganic chemical compounds as its source of energy. Almost all life on Earth relies directly or indirectly on primary production. The organisms responsible for primary production are known as primary producers or autotrophs, and form the base of the food chain. In terrestrial ecoregions, these are mainly plants, while in aquatic ecoregions algae predominate in this role. Ecologists distinguish primary production as either net or gross, the former accounting for losses to processes such as cellular respiration, the latter not. Overview Primary production is the production of chemical energy in organic compounds by living organisms. The main source of this energy is sunlight but a minute fraction of primary production is driven by lithotrophic organisms using the chemical energy of inorganic molecules.Regardless of its source, this energy is used to synthesize complex organic molecules from simpler inorganic compounds such as carbon dioxide () and water (H2O). The following two equations are simplified representations of photosynthesis (top) and (one form of) chemosynthesis (bottom): + H2O + light → CH2O + O2 + O2 + 4 H2S → CH2O + 4 S + 3 H2O In both cases, the end point is a polymer of reduced carbohydrate, (CH2O)n, typically molecules such as glucose or other sugars. These relatively simple molecules may be then used to further synthesise more complicated molecules, including proteins, complex carbohydrates, lipids, and nucleic acids, or be respired to perform work. Consumption of primary producers by heterotrophic organisms, such as animals, then transfers these organic molecules (and the energy stored within them) up the food web, fueling all of the Earth's living systems. Gross primary production and net primary production Gross primary production (GPP) is the amount of chemical energy, typically expressed as carbon biomass, that primary producers create in a given length of time. Some fraction of this fixed energy is used by primary producers for cellular respiration and maintenance of existing tissues (i.e., "growth respiration" and "maintenance respiration"). The remaining fixed energy (i.e., mass of photosynthate) is referred to as net primary production (NPP). NPP = GPP - respiration [by plants] Net primary production is the rate at which all the autotrophs in an ecosystem produce net useful chemical energy. Net primary production is available to be directed toward growth and reproduction of primary producers. As such it is available for consumption by herbivores. Both gross and net primary production are typically expressed in units of mass per unit area per unit time interval. In terrestrial ecosystems, mass of carbon per unit area per year (g C m−2 yr−1) is most often used as the unit of measurement. Note that a distinction is sometimes drawn between "production" and "productivity", with the former the quantity of material produced (g C m−2), the latter the rate at which it is produced (g C m−2 yr−1), but these terms are more typically used interchangeably. Terrestrial production On the land, almost all primary production is now performed by vascular plants, with a small fraction coming from algae and non-vascular plants such as mosses and liverworts. Before the evolution of vascular plants, non-vascular plants likely played a more significant role. Primary production on land is a function of many factors, but principally local hydrology and temperature (the latter covaries to an extent with light, specifically photosynthetically active radiation (PAR), the source of energy for photosynthesis). While plants cover much of the Earth's surface, they are strongly curtailed wherever temperatures are too extreme or where necessary plant resources (principally water and PAR) are limiting, such as deserts or polar regions. Water is "consumed" in plants by the processes of photosynthesis (see above) and transpiration. The latter process (which is responsible for about 90% of water use) is driven by the evaporation of water from the leaves of plants. Transpiration allows plants to transport water and mineral nutrients from the soil to growth regions, and also cools the plant. Diffusion of water vapour out of a leaf, the force that drives transpiration, is regulated by structures known as stomata. These structures also regulate the diffusion of carbon dioxide from the atmosphere into the leaf, such that decreasing water loss (by partially closing stomata) also decreases carbon dioxide gain. Certain plants use alternative forms of photosynthesis, called Crassulacean acid metabolism (CAM) and C4. These employ physiological and anatomical adaptations to increase water-use efficiency and allow increased primary production to take place under conditions that would normally limit carbon fixation by C3 plants (the majority of plant species). As shown in the animation, the boreal forests of Canada and Russia experience high productivity in June and July and then a slow decline through fall and winter. Year-round, tropical forests in South America, Africa, Southeast Asia, and Indonesia have high productivity, not surprising with the abundant sunlight, warmth, and rainfall. However, even in the tropics, there are variations in productivity over the course of the year. For example, the Amazon basin exhibits especially high productivity from roughly August through October - the period of the area's dry season. Because the trees have access to a plentiful supply of ground water that builds up in the rainy season, they grow better when the rainy skies clear and allow more sunlight to reach the forest. Oceanic production In a reversal of the pattern on land, in the oceans, almost all photosynthesis is performed by algae, with a small fraction contributed by vascular plants and other groups. Algae encompass a diverse range of organisms, ranging from single floating cells to attached seaweeds. They include photoautotrophs from a variety of groups. Eubacteria are important photosynthetizers in both oceanic and terrestrial ecosystems, and while some archaea are phototrophic, none are known to utilise oxygen-evolving photosynthesis. A number of eukaryotes are significant contributors to primary production in the ocean, including green algae, brown algae and red algae, and a diverse group of unicellular groups. Vascular plants are also represented in the ocean by groups such as the seagrasses. Unlike terrestrial ecosystems, the majority of primary production in the ocean is performed by free-living microscopic organisms called phytoplankton. Larger autotrophs, such as the seagrasses and macroalgae (seaweeds) are generally confined to the littoral zone and adjacent shallow waters, where they can attach to the underlying substrate but still be within the photic zone. There are exceptions, such as Sargassum, but the vast majority of free-floating production takes place within microscopic organisms. The factors limiting primary production in the ocean are also very different from those on land. The availability of water, obviously, is not an issue (though its salinity can be). Similarly, temperature, while affecting metabolic rates (see Q10), ranges less widely in the ocean than on land because the heat capacity of seawater buffers temperature changes, and the formation of sea ice insulates it at lower temperatures. However, the availability of light, the source of energy for photosynthesis, and mineral nutrients, the building blocks for new growth, play crucial roles in regulating primary production in the ocean. Available Earth System Models suggest that ongoing ocean bio-geochemical changes could trigger reductions in ocean NPP between 3% and 10% of current values depending on the emissions scenario. Light The sunlit zone of the ocean is called the photic zone (or euphotic zone). This is a relatively thin layer (10–100 m) near the ocean's surface where there is sufficient light for photosynthesis to occur. For practical purposes, the thickness of the photic zone is typically defined by the depth at which light reaches 1% of its surface value. Light is attenuated down the water column by its absorption or scattering by the water itself, and by dissolved or particulate material within it (including phytoplankton). Net photosynthesis in the water column is determined by the interaction between the photic zone and the mixed layer. Turbulent mixing by wind energy at the ocean's surface homogenises the water column vertically until the turbulence dissipates (creating the aforementioned mixed layer). The deeper the mixed layer, the lower the average amount of light intercepted by phytoplankton within it. The mixed layer can vary from being shallower than the photic zone, to being much deeper than the photic zone. When it is much deeper than the photic zone, this results in phytoplankton spending too much time in the dark for net growth to occur. The maximum depth of the mixed layer in which net growth can occur is called the critical depth. As long as there are adequate nutrients available, net primary production occurs whenever the mixed layer is shallower than the critical depth. Both the magnitude of wind mixing and the availability of light at the ocean's surface are affected across a range of space- and time-scales. The most characteristic of these is the seasonal cycle (caused by the consequences of the Earth's axial tilt), although wind magnitudes additionally have strong spatial components. Consequently, primary production in temperate regions such as the North Atlantic is highly seasonal, varying with both incident light at the water's surface (reduced in winter) and the degree of mixing (increased in winter). In tropical regions, such as the gyres in the middle of the major basins, light may only vary slightly across the year, and mixing may only occur episodically, such as during large storms or hurricanes. Nutrients Mixing also plays an important role in the limitation of primary production by nutrients. Inorganic nutrients, such as nitrate, phosphate and silicic acid are necessary for phytoplankton to synthesise their cells and cellular machinery. Because of gravitational sinking of particulate material (such as plankton, dead or fecal material), nutrients are constantly lost from the photic zone, and are only replenished by mixing or upwelling of deeper water. This is exacerbated where summertime solar heating and reduced winds increases vertical stratification and leads to a strong thermocline, since this makes it more difficult for wind mixing to entrain deeper water. Consequently, between mixing events, primary production (and the resulting processes that leads to sinking particulate material) constantly acts to consume nutrients in the mixed layer, and in many regions this leads to nutrient exhaustion and decreased mixed layer production in the summer (even in the presence of abundant light). However, as long as the photic zone is deep enough, primary production may continue below the mixed layer where light-limited growth rates mean that nutrients are often more abundant. Iron Another factor relatively recently discovered to play a significant role in oceanic primary production is the micronutrient iron. This is used as a cofactor in enzymes involved in processes such as nitrate reduction and nitrogen fixation. A major source of iron to the oceans is dust from the Earth's deserts, picked up and delivered by the wind as aeolian dust. In regions of the ocean that are distant from deserts or that are not reached by dust-carrying winds (for example, the Southern and North Pacific oceans), the lack of iron can severely limit the amount of primary production that can occur. These areas are sometimes known as HNLC (High-Nutrient, Low-Chlorophyll) regions, because the scarcity of iron both limits phytoplankton growth and leaves a surplus of other nutrients. Some scientists have suggested introducing iron to these areas as a means of increasing primary productivity and sequestering carbon dioxide from the atmosphere. Measurement The methods for measurement of primary production vary depending on whether gross vs net production is the desired measure, and whether terrestrial or aquatic systems are the focus. Gross production is almost always harder to measure than net, because of respiration, which is a continuous and ongoing process that consumes some of the products of primary production (i.e. sugars) before they can be accurately measured. Also, terrestrial ecosystems are generally more difficult because a substantial proportion of total productivity is shunted to below-ground organs and tissues, where it is logistically difficult to measure. Shallow water aquatic systems can also face this problem. Scale also greatly affects measurement techniques. The rate of carbon assimilation in plant tissues, organs, whole plants, or plankton samples can be quantified by biochemically based techniques, but these techniques are decidedly inappropriate for large scale terrestrial field situations. There, net primary production is almost always the desired variable, and estimation techniques involve various methods of estimating dry-weight biomass changes over time. Biomass estimates are often converted to an energy measure, such as kilocalories, by an empirically determined conversion factor. Terrestrial In terrestrial ecosystems, researchers generally measure net primary production (NPP). Although its definition is straightforward, field measurements used to estimate productivity vary according to investigator and biome. Field estimates rarely account for below ground productivity, herbivory, turnover, litterfall, volatile organic compounds, root exudates, and allocation to symbiotic microorganisms. Biomass based NPP estimates result in underestimation of NPP due to incomplete accounting of these components. However, many field measurements correlate well to NPP. There are a number of comprehensive reviews of the field methods used to estimate NPP. Estimates of ecosystem respiration, the total carbon dioxide produced by the ecosystem, can also be made with gas flux measurements. The major unaccounted pool is belowground productivity, especially production and turnover of roots. Belowground components of NPP are difficult to measure. BNPP (below-ground NPP) is often estimated based on a ratio of ANPP:BNPP (above-ground NPP:below-ground NPP) rather than direct measurements. Gross primary production can be estimated from measurements of net ecosystem exchange (NEE) of carbon dioxide made by the eddy covariance technique. During night, this technique measures all components of ecosystem respiration. This respiration is scaled to day-time values and further subtracted from NEE. Grasslands Most frequently, peak standing biomass is assumed to measure NPP. In systems with persistent standing litter, live biomass is commonly reported. Measures of peak biomass are more reliable if the system is predominantly annuals. However, perennial measurements could be reliable if there were a synchronous phenology driven by a strong seasonal climate. These methods may underestimate ANPP in grasslands by as much as 2 (temperate) to 4 (tropical) fold. Repeated measures of standing live and dead biomass provide more accurate estimates of all grasslands, particularly those with large turnover, rapid decomposition, and interspecific variation in timing of peak biomass. Wetland productivity (marshes and fens) is similarly measured. In Europe, annual mowing makes the annual biomass increment of wetlands evident. Forests Methods used to measure forest productivity are more diverse than those of grasslands. Biomass increment based on stand specific allometry plus litterfall is considered a suitable although incomplete accounting of above-ground net primary production (ANPP). Field measurements used as a proxy for ANPP include annual litterfall, diameter or basal area increment (DBH or BAI), and volume increment. Aquatic In aquatic systems, primary production is typically measured using one of six main techniques: variations in oxygen concentration within a sealed bottle (developed by Gaarder and Gran in 1927) incorporation of inorganic carbon-14 (14C in the form of sodium bicarbonate) into organic matter Stable isotopes of Oxygen (16O, 18O and 17O) fluorescence kinetics (technique still a research topic) Stable isotopes of Carbon (12C and 13C) Oxygen/Argon Ratios The technique developed by Gaarder and Gran uses variations in the concentration of oxygen under different experimental conditions to infer gross primary production. Typically, three identical transparent vessels are filled with sample water and stoppered. The first is analysed immediately and used to determine the initial oxygen concentration; usually this is done by performing a Winkler titration. The other two vessels are incubated, one each in under light and darkened. After a fixed period of time, the experiment ends, and the oxygen concentration in both vessels is measured. As photosynthesis has not taken place in the dark vessel, it provides a measure of ecosystem respiration. The light vessel permits both photosynthesis and respiration, so provides a measure of net photosynthesis (i.e. oxygen production via photosynthesis subtract oxygen consumption by respiration). Gross primary production is then obtained by adding oxygen consumption in the dark vessel to net oxygen production in the light vessel. The technique of using 14C incorporation (added as labelled Na2CO3) to infer primary production is most commonly used today because it is sensitive, and can be used in all ocean environments. As 14C is radioactive (via beta decay), it is relatively straightforward to measure its incorporation in organic material using devices such as scintillation counters. Depending upon the incubation time chosen, net or gross primary production can be estimated. Gross primary production is best estimated using relatively short incubation times (1 hour or less), since the loss of incorporated 14C (by respiration and organic material excretion / exudation) will be more limited. Net primary production is the fraction of gross production remaining after these loss processes have consumed some of the fixed carbon. Loss processes can range between 10 and 60% of incorporated 14C according to the incubation period, ambient environmental conditions (especially temperature) and the experimental species used. Aside from those caused by the physiology of the experimental subject itself, potential losses due to the activity of consumers also need to be considered. This is particularly true in experiments making use of natural assemblages of microscopic autotrophs, where it is not possible to isolate them from their consumers. The methods based on stable isotopes and O2/Ar ratios have the advantage of providing estimates of respiration rates in the light without the need of incubations in the dark. Among them, the method of the triple oxygen isotopes and O2/Ar have the additional advantage of not needing incubations in closed containers and O2/Ar can even be measured continuously at sea using equilibrator inlet mass spectrometry (EIMS) or a membrane inlet mass spectrometry (MIMS). However, if results relevant to the carbon cycle are desired, it is probably better to rely on methods based on carbon (and not oxygen) isotopes. It is important to notice that the method based on carbon stable isotopes is not simply an adaptation of the classic 14C method, but an entirely different approach that does not suffer from the problem of lack of account of carbon recycling during photosynthesis. Global As primary production in the biosphere is an important part of the carbon cycle, estimating it at the global scale is important in Earth system science. However, quantifying primary production at this scale is difficult because of the range of habitats on Earth, and because of the impact of weather events (availability of sunlight, water) on its variability. Using satellite-derived estimates of the Normalized Difference Vegetation Index (NDVI) for terrestrial habitats and sea-surface chlorophyll for the oceans, it is estimated that the total (photoautotrophic) primary production for the Earth was 104.9 petagrams of carbon per year (Pg C yr−1; equivalent to the non-SI Gt C yr−1). Of this, 56.4 Pg C yr−1 (53.8%), was the product of terrestrial organisms, while the remaining 48.5 Pg C yr−1, was accounted for by oceanic production. Scaling ecosystem-level GPP estimations based on eddy covariance measurements of net ecosystem exchange (see above) to regional and global values using spatial details of different predictor variables, such as climate variables and remotely sensed fAPAR or LAI led to a terrestrial gross primary production of 123±8 Gt carbon (NOT carbon dioxide) per year during 1998-2005 In areal terms, it was estimated that land production was approximately 426 g C m−2 yr−1 (excluding areas with permanent ice cover), while that for the oceans was 140 g C m−2 yr−1. Another significant difference between the land and the oceans lies in their standing stocks - while accounting for almost half of total production, oceanic autotrophs only account for about 0.2% of the total biomass. Present and Past Estimates Present day primary productivity can be estimated through a variety of methodologies including ship-board measurements, satellites and terrestrial observatories. Historical estimates have relied on biogeochemical models and geochemical proxies. One example is using barium, where barite concentrations in marine sediments rise in line with carbon export production at the surface. Another example is using the triple oxygen isotopes of sulfate. Together these records suggest large shifts in primary production throughout Earth's past with notable rises associated with Earth's Great Oxidation Event (approximately 2.4 to 2.0 billion years ago) and the Neoproterozoic (approximately 1.0 to 0.54 billion years ago). Human impact and appropriation Human societies are part of the Earth's NPP cycle but disproportionately influence it. In 1996, Josep Garí designed a new indicator of sustainable development based precisely on the estimation of the human appropriation of NPP: he coined it "HANPP" (Human Appropriation of Net Primary Production) and introduced it at the inaugural conference of the European Society for Ecological Economics. HANPP has since been further developed and widely applied in research on ecological economics and in policy analysis for sustainability. HANPP represents a proxy of the human impact on nature and can be applied to different geographical and global scales. The extensive degree of human use of the Planet's resources, mostly via land use, results in various levels of impact on actual NPP (NPPact). Although in some regions, such as the Nile valley, irrigation has resulted in a considerable increase in primary production, in most of the Planet, there is a notable trend of NPP reduction due to land changes (ΔNPPLC) of 9.6% across global land-mass. In addition to this, end consumption by people raises the total HANPP to 23.8% of potential vegetation (NPP0). It is estimated that, in 2000, 34% of the Earth's ice-free land area (12% cropland; 22% pasture) was devoted to human agriculture. This disproportionate amount reduces the energy available to other species, having a marked impact on biodiversity, flows of carbon, water, and energy, and ecosystem services,. Scientists have questioned how large this fraction can be before these services break down. Reductions in NPP are also expected in the ocean as a result of ongoing climate change, potentially impacting marine ecosystems (~10% of global biodiversity) and goods and services (1-5% of global total) that the oceans provide.	Biology and health sciences	Ecology	Biology
208318	https://en.wikipedia.org/wiki/Cormorant	Cormorant	Phalacrocoracidae is a family of approximately 42 species of aquatic birds commonly known as cormorants and shags. Several different classifications of the family have been proposed, but in 2021 the International Ornithologists' Union (IOU) adopted a consensus taxonomy of seven genera. The great cormorant (Phalacrocorax carbo) and the common shag (Gulosus aristotelis) are the only two species of the family commonly encountered in Britain and Ireland and "cormorant" and "shag" appellations have been later assigned to different species in the family somewhat haphazardly. Cormorants and shags are medium-to-large birds, with body weight in the range of and wing span of . The majority of species have dark feathers. The bill is long, thin and hooked. Their feet have webbing between all four toes. All species are fish-eaters, catching the prey by diving from the surface. They are excellent divers, and under water they propel themselves with their feet with help from their wings; some cormorant species have been found to dive as deep as . They have relatively short wings due to their need for economical movement underwater, and consequently have among the highest flight costs of any flying bird. Cormorants nest in colonies around the shore, on trees, islets or cliffs. They are coastal rather than oceanic birds, and some have colonised inland waters. The original ancestor of cormorants seems to have been a fresh-water bird. They range around the world, except for the central Pacific islands. Names "Cormorant" is a contraction probably derived from Latin corvus marinus, "sea raven". Cormoran is the Cornish name of the sea giant in the tale of Jack the Giant Killer. Indeed, "sea raven" or analogous terms were the usual terms for cormorants in Germanic languages until after the Middle Ages. The French explorer André Thévet commented in 1558: "the beak [is] similar to that of a cormorant or other corvid", which demonstrates that the erroneous belief that the birds were related to ravens lasted at least to the 16th century. No consistent distinction exists between cormorants and shags. The names "cormorant" and "shag" were originally the common names of the two species of the family found in Great Britain Phalacrocorax carbo (now referred to by ornithologists as the great cormorant) and Gulosus aristotelis (the European shag). "Shag" refers to the bird's crest, which the British forms of the great cormorant lack. As other species were encountered by English-speaking sailors and explorers elsewhere in the world, some were called cormorants and some shags, sometimes depending on whether they had crests or not. Sometimes the same species is called a cormorant in one part of the world and a shag in another; for example, all species in the family which occur in New Zealand are known locally as shags, including four non-endemic species known as cormorant elsewhere in their range. Van Tets (1976) proposed to divide the family into two genera and attach the name "cormorant" to one and "shag" to the other, but this nomenclature has not been widely adopted. Description Cormorants and shags are medium-to-large seabirds. They range in size from the pygmy cormorant (Microcarbo pygmaeus), at as little as and , to the flightless cormorant (Nannopterum harrisi), at a maximum size and . The recently extinct spectacled cormorant (Urile perspicillatus) was rather larger, at an average size of . The majority, including nearly all Northern Hemisphere species, have mainly dark plumage, but some Southern Hemisphere species are black and white, and a few (e.g. the spotted shag of New Zealand) are quite colourful. Many species have areas of coloured skin on the face (the lores and the gular skin) which can be bright blue, orange, red or yellow, typically becoming more brightly coloured in the breeding season. The bill is long, thin, and sharply hooked. Their feet have webbing between all four toes, as in their relatives. Habitat Habitat varies from species to species: some are restricted to seacoasts, while others occur in both coastal and inland waters to varying degrees. They range around the world, except for the central Pacific islands. Behaviour All cormorants and shags are fish-eaters, dining on small eels, fish, and even water snakes. They dive from the surface, though many species make a characteristic half-jump as they dive, presumably to give themselves a more streamlined entry into the water. Under water they propel themselves with their feet, though some also propel themselves with their wings (see the picture, commentary, and existing reference video). Imperial shags fitted with miniaturized video recorders have been filmed diving to depths of as much as to forage on the sea floor. After fishing, cormorants go ashore, and are frequently seen holding their wings out in the sun. All cormorants have preen gland secretions that are used ostensibly to keep the feathers waterproof. Some sources state that cormorants have waterproof feathers while others say that they have water-permeable feathers. Still others suggest that the outer plumage absorbs water but does not permit it to penetrate the layer of air next to the skin. The wing drying action is seen even in the flightless cormorant but not in the Antarctic shags or red-legged cormorants. Alternate functions suggested for the spread-wing posture include that it aids thermoregulation or digestion, balances the bird, or indicates presence of fish. A detailed study of the great cormorant concluded there is little doubt that is serves to dry the plumage. Cormorants are colonial nesters, using trees, rocky islets, or cliffs. The eggs are a chalky-blue colour. There is usually one brood a year. Parents regurgitate food to feed their young. Taxonomy The genus Phalacrocorax, from which the family name Phalacrocoracidae is derived, is Latinised from Ancient Greek phalakros "bald" and korax "raven". This is thought to refer to the creamy white patch on the cheeks of adult great cormorants, or the ornamental white head plumes prominent in Mediterranean birds of this species, but is certainly not a unifying characteristic of cormorants. The cormorant family was traditionally placed within the Pelecaniformes or, in the Sibley–Ahlquist taxonomy of the 1990s, the expanded Ciconiiformes. Pelecaniformes in the traditional sense—all waterbird groups with totipalmate foot webbing—are not a monophyletic group, even after the removal of the distantly-related tropicbirds. Their relationships and delimitation – apart from being part of a "higher waterfowl" clade which is similar but not identical to Sibley and Ahlquist's "pan-Ciconiiformes" – remain mostly unresolved. Notwithstanding, all evidence agrees that the cormorants and shags are closer to the darters and Sulidae (gannets and boobies), and perhaps the pelicans or even penguins, than to all other living birds. In recent years, three preferred treatments of the cormorant family have emerged: either to leave all living cormorants in a single genus, Phalacrocorax, or to split off a few species such as the imperial shag complex (in Leucocarbo) and perhaps the flightless cormorant. Alternatively, the genus may be disassembled altogether and in the most extreme case be reduced to the great, white-breasted and Japanese cormorants. In 2014, a landmark study proposed a 7 genera treatment, which was adopted by the IUCN Red List and BirdLife International, and later by the IOC in 2021, standardizing it. The cormorants and the darters have a unique bone on the back of the top of the skull known as the os nuchale or occipital style which was called a xiphoid process in early literature. This bony projection provides anchorage for the muscles that increase the force with which the lower mandible is closed. This bone and the highly developed muscles over it, the M. adductor mandibulae caput nuchale, are unique to the families Phalacrocoracidae and Anhingidae. Several evolutionary groups are still recognizable. However, combining the available evidence suggests that there has also been a great deal of convergent evolution; for example the cliff shags are a convergent paraphyletic group. The proposed division into Phalacrocorax sensu stricto (or subfamily "Phalacrocoracinae") cormorants and Leucocarbo sensu lato (or "Leucocarboninae") shags does have some degree of merit. The resolution provided by the mtDNA 12S rRNA and ATPase subunits six and eight sequence data is not sufficient to properly resolve several groups to satisfaction; in addition, many species remain unsampled, the fossil record has not been integrated in the data, and the effects of hybridisation – known in some Pacific species especially – on the DNA sequence data are unstudied. A multigene molecular phylogenetic study published in 2014 provided a genus-level phylogeny of the family. List of genera As per the IOU, the IUCN Red List and BirdLife International, the family contains 7 genera: Prior to 2021, the IOU (or formerly the IOC) classified all these species in just three genera: Microcarbo, Leucocarbo, and a broad Phalacrocorax containing all remaining species; however, this treatment rendered Phalacrocorax deeply paraphyletic with respect to Leucocarbo. Other authorities, such as the Clements Checklist, formerly recognised only Microcarbo as a separate genus from Phalacrocorax. For details, see the article "List of cormorant species". Evolution and fossil record The details of the evolution of the cormorants are mostly unknown. Even the technique of using the distribution and relationships of a species to figure out where it came from, biogeography, usually very informative, does not give very specific data for this probably rather ancient and widespread group. However, the closest living relatives of the cormorants and shags are the other families of the suborder Sulae—darters and gannets and boobies—which have a primarily Gondwanan distribution. Hence, at least the modern diversity of Sulae probably originated in the southern hemisphere. While the Leucocarbonines are almost certainly of southern Pacific origin—possibly even the Antarctic which, at the time when cormorants evolved, was not yet ice-covered—all that can be said about the Phalacrocoracines is that they are most diverse in the regions bordering the Indian Ocean, but generally occur over a large area. Similarly, the origin of the family is shrouded in uncertainties. Some Late Cretaceous fossils have been proposed to belong with the Phalacrocoracidae: A scapula from the Campanian-Maastrichtian boundary, about 70 mya (million years ago), was found in the Nemegt Formation in Mongolia; it is now in the PIN collection. It is from a bird roughly the size of a spectacled cormorant, and quite similar to the corresponding bone in Phalacrocorax. A Maastrichtian (Late Cretaceous, c. 66 mya) right femur, AMNH FR 25272 from the Lance Formation near Lance Creek, Wyoming, is sometimes suggested to be the second-oldest record of the Phalacrocoracidae; this was from a rather smaller bird, about the size of a long-tailed cormorant. However, cormorants likely originated much later, and these are likely misidentifications. As the Early Oligocene "Sula" ronzoni cannot be assigned to any of the sulid families—cormorants and shags, darters, and gannets and boobies—with certainty, the best interpretation is that the Phalacrocoracidae diverged from their closest ancestors in the Early Oligocene, perhaps some 30 million years ago, and that the Cretaceous fossils represent ancestral sulids, "pelecaniforms" or "higher waterbirds"; at least the last lineage is generally believed to have been already distinct and undergoing evolutionary radiation at the end of the Cretaceous. What can be said with near certainty is that AMNH FR 25272 is from a diving bird that used its feet for underwater locomotion; as this is liable to result in some degree of convergent evolution and the bone is missing indisputable neornithine features, it is not entirely certain that the bone is correctly referred to this group. Phylogenetic evidence indicates that the cormorants diverged from their closest relatives, the darters, during the Late Oligocene, indicating that most of the claims of Cretaceous or early Paleogene cormorant occurrences are likely misidentifications. During the late Paleogene, when the family presumably originated, much of Eurasia was covered by shallow seas, as the Indian Plate finally attached to the mainland. Lacking a detailed study, it may well be that the first "modern" cormorants were small species from eastern, south-eastern or southern Asia, possibly living in freshwater habitat, that dispersed due to tectonic events. Such a scenario would account for the present-day distribution of cormorants and shags and is not contradicted by the fossil record; as remarked above, a thorough review of the problem is not yet available. Even when Phalacrocorax was used to unite all living species, two distinct genera of prehistoric cormorants became widely accepted today: Limicorallus (Indricotherium middle Oligocene of Chelkar-Teniz, Kazakhstan) Nectornis (Late Oligocene/Early Miocene of Central Europe – Middle Miocene of Bes-Konak, Turkey) – includes Oligocorax miocaenus The proposed genus Oligocorax appears to be paraphyletic – the European species have been separated in Nectornis, and the North American ones placed in the expanded Phalacrocorax; the latter might just as well be included in Nannopterum. A Late Oligocene fossil cormorant foot from Enspel, Germany, sometimes placed in Oligocorax, would then be referable to Nectornis if it proves not to be too distinct. Limicorallus, meanwhile, was initially believed to be a rail or a dabbling duck by some. There are also undescribed remains of apparent cormorants from the Quercy Phosphorites of Quercy (France), dating to some time between the Late Eocene and the mid-Oligocene. All these early European species might belong to the basal group of "microcormorants", as they conform with them in size and seem to have inhabited the same habitat: subtropical coastal or inland waters. While this need not be more than convergence, the phylogeny of the modern (sub)genus Microcarbo – namely, whether the Western Eurasian M. pygmaeus is a basal or highly derived member of its clade – is still not well understood at all as of 2022. Some other Paleogene remains are sometimes assigned to the Phalacrocoracidae, but these birds seem rather intermediate between cormorants and darters (and lack clear autapomorphies of either). Thus, they may be quite basal members of the Palacrocoracoidea. The taxa in question are: Piscator (Late Eocene of England) "Pelecaniformes" gen. et sp. indet. (Jebel Qatrani Early Oligocene of Fayum, Egypt) – similar to Piscator? Borvocarbo (Late Oligocene of C Europe) The supposed Late Pliocene/Early Pleistocene "Valenticarbo" is a nomen dubium and given its recent age probably not a separate genus. The remaining fossil species are not usually placed in a modern phylogenetic framework. While the numerous western US species are most likely prehistoric representatives of the coastal Urile or inland Nannopterum, the European fossils pose much more of a problem due to the singular common shag being intermediate in size between the other two European cormorant lineages, and as of 2022 still of mysterious ancestry; notably, a presumably lost collection of Late Miocene fossils from the Odesa region may have contained remains of all three (sub)genera inhabiting Europe today. Similarly, the Plio-Pleistocene fossils from Florida have been allied with Nannopterum and even Urile, but may conceivably be Phalacrocorax; they are in serious need of revision since it is not even clear how many species are involved. Provisionally, the fossil species are thus all placed in Phalacrocorax here: Phalacrocorax marinavis (Oligocene – Early Miocene of Oregon, US) – formerly Oligocorax; Urile or Nannopterum? Phalacrocorax littoralis (Late Oligocene/Early Miocene of St-Gérand-le-Puy, France) – formerly Oligocorax; Nectornis? Phalacrocorax intermedius (Early – Middle Miocene of C Europe) – includes P. praecarbo, Ardea/P. brunhuberi and Botaurites avitus; Microcarbo, Phalacrocorax or Gulosus? Phalacrocorax macropus (Early Miocene – Pliocene of north-west US) – Urile or Nannopterum? Phalacrocorax ibericus (Late Miocene of Valles de Fuentiduena, Spain) – Microcarbo, Phalacrocorax or Gulosus? Phalacrocorax lautus (Late Miocene of Golboçica, Moldavia) – Microcarbo, Phalacrocorax or Gulosus? Phalacrocorax serdicensis (Late Miocene of Hrabarsko, Bulgaria); Microcarbo, Phalacrocorax or Gulosus? Phalacrocorax sp(p). (Late Miocene of Odesa region, Ukraine) – up to 4 species, one of which is probably P. longipes; Microcarbo, Phalacrocorax and/or Gulosus? Phalacrocorax femoralis (Modelo Late Miocene/Early Pliocene of WC North America) – formerly Miocorax; Nannopterum? Phalacrocorax sp. (Late Miocene/Early Pliocene of Lee Creek Mine, US) – Nannopterum or Phalacrocorax? Phalacrocorax sp. 1 (Late Miocene/Early Pliocene of WC South America) – probably Leucocarbo Phalacrocorax sp. 2 (Pisco Late Miocene/Early Pliocene of SW Peru) – Poikilocarbo or Leucocarbo? Phalacrocorax longipes (Late Miocene – Early Pliocene of Ukraine) – formerly Pliocarbo; Microcarbo, Phalacrocorax or Gulosus? Phalacrocorax goletensis (Early Pliocene – Early Pleistocene of Mexico) – Urile or Nannopterum, perhaps Poikilocarbo or Leuocarbo Phalacrocorax wetmorei (Bone Valley Early Pliocene of Florida) – Nannopterum or Phalacrocorax? Phalacrocorax sp. (Bone Valley Early Pliocene of Polk County, Florida, US) – Nannopterum or Phalacrocorax? Phalacrocorax leptopus (Juntura Early/Middle Pliocene of Juntura, Malheur County, Oregon, US) – Nannopterum? Phalacrocorax reliquus (Middle Pliocene of Mongolia) – Microcarbo, Phalacrocorax or Gulosus? Phalacrocorax idahensis (Middle Pliocene – Pleistocene of Idaho, US, and possibly Florida) – Nannopterum? Phalacrocorax destefanii (Late Pliocene of Italy) – formerly Paracorax; Microcarbo, Phalacrocorax or Gulosus? Phalacrocorax filyawi (Pinecrest Late Pliocene of Florida, US) – may be P. idahensis; Nannopterum or Phalacrocorax, perhaps Urile? Phalacrocorax kennelli (San Diego Late Pliocene of California, US) – Urile or Nannopterum? Phalacrocorax kumeyaay (San Diego Late Pliocene of California, US) – Urile or Nannopterum? Phalacrocorax macer (Late Pliocene of Idaho, US) – Nannopterum? Phalacrocorax mongoliensis (Late Pliocene of W Mongolia) – Microcarbo, Phalacrocorax or Gulosus? Phalacrocorax sp. (La Portada Late Pliocene of N Chile) – may be same as Late Miocene/Early Pliocene "Phalacrocorax sp. 2"; Poikilocarbo or Leucocarbo? Phalacrocorax rogersi (Late Pliocene – Early Pleistocene of California, US) – Urile or Nannopterum? Phalacrocorax chapalensis (Late Pliocene/Early Pleistocene of Jalisco, Mexico) – Urile or Nannopterum, perhaps Poikilocarbo or Leucocarbo? Phalacrocorax gregorii (Late Pleistocene of Australia) – possibly not a valid species; Microcarbo, Phalacrocorax or Leucocarbo? Phalacrocorax vetustus (Late Pleistocene of Australia) – formerly Australocorax, possibly not a valid species; Microcarbo, Phalacrocorax or Leucocarbo? Phalacrocorax sp. (Sarasota County, Florida, US) – may be P. filawyi/idahensis; Nannopterum or Phalacrocorax? The former "Phalacrocorax" (or "Oligocorax") mediterraneus is now considered to belong to the bathornithid Paracrax antiqua. "P." subvolans was actually a darter (Anhinga). In human culture Cormorant culling Cormorant fishing Humans have used cormorants' fishing skills in various places in the world. Archaeological evidence suggests that cormorant fishing was practised in Ancient Egypt, Peru, Korea and India, but the strongest tradition has remained in China and Japan, where it reached commercial-scale level in some areas. In Japan, cormorant fishing is called and is performed by a fisherman known as an usho. Traditional forms of ukai can be seen on the Nagara River in the city of Gifu, Gifu Prefecture, where cormorant fishing has continued uninterrupted for 1300 years, or in the city of Inuyama, Aichi. In Guilin, Guangxi, cormorants are famous for fishing on the shallow Li River. In Gifu, the Japanese cormorant (P. capillatus) is used; Chinese fishermen often employ great cormorants (P. carbo). In Europe, a similar practice was also used on Doiran Lake in the region of Macedonia. James VI and I appointed a keeper of cormorants, John Wood, and built ponds at Westminster to train the birds to fish. In a common technique, a snare is tied near the base of the bird's throat, which allows the bird only to swallow small fish. When the bird captures and tries to swallow a large fish, the fish is caught in the bird's throat. When the bird returns to the fisherman's raft, the fisherman helps the bird to remove the fish from its throat. The method is not as common today, since more efficient methods of catching fish have been developed, but is still practised as a cultural tradition. In Japan, environmental changes threaten traditional ukai because of reduced numbers of the ayu river fish that cormorants are used to catch. In folklore, literature, and art Cormorants feature in heraldry and medieval ornamentation, usually in their "wing-drying" pose, which was seen as representing the Christian cross, and symbolizing nobility and sacrifice. For John Milton in Paradise Lost, the cormorant symbolizes greed: perched atop the Tree of Life, Satan took the form of a cormorant as he spied on Adam and Eve during his first intrusion into Eden. In some Scandinavian areas, they are considered good omen; in particular, in Norwegian tradition spirits of those lost at sea come to visit their loved ones disguised as cormorants. For example, the Norwegian municipalities of Røst, Loppa and Skjervøy have cormorants in their coat of arms. The symbolic liver bird of Liverpool is commonly thought to be a cross between an eagle and a cormorant. In Homer's epic poem The Odyssey, Odysseus (Ulysses) is saved by a compassionate sea nymph who takes the form of a cormorant. In 1853, a woman wearing a dress made of cormorant feathers was found on San Nicolas Island, off the southern coast of California. She had sewn the feather dress together using whale sinews. She is known as the Lone Woman of San Nicolas and was later baptised "Juana Maria" (her original name is lost). The woman had lived alone on the island for 18 years before being rescued. When removed from San Nicolas, she brought with her a green cormorant dress she made; this dress is reported to have been removed to the Vatican. Her story, which includes the feather dress, was fictionalized in the children's novel Island of the Blue Dolphins. The bird has inspired numerous writers, including Amy Clampitt, who wrote a poem called "The Cormorant in its Element". The species she described may have been the pelagic cormorant, which is the only species in the temperate U.S. with the "slim head ... vermilion-strapped" and "big black feet" that she mentions. A cormorant representing Blanche Ingram appears in the first of the fictional paintings by Jane in Charlotte Brontë's novel Jane Eyre: One gleam of light lifted into relief a half-submerged mast, on which sat a cormorant, dark and large, with wings flecked with foam; its beak held a gold bracelet, set with gems, that I had touched with as brilliant tints as my palette could yield, and as glittering distinctness as my pencil could impart. In the Sherlock Holmes story "The Adventure of the Veiled Lodger", Dr. Watson warns that if there are further attempts to get at and destroy his private notes regarding his time with Holmes, "the whole story concerning the politician, the lighthouse, and the trained cormorant will be given to the public. There is at least one reader who will understand." A cormorant is humorously mentioned as having had linseed oil rubbed into it by a wayward pupil during the "Growth and Learning" segment of the 1983 Monty Python movie Monty Python's The Meaning of Life. The cormorant served as the hood ornament for the Packard automobile brand. Cormorants (and books about them written by a fictional ornithologist) are a recurring fascination of the protagonist in Jesse Ball's 2018 novel Census. The Pokémon Cramorant, featured in the 8th generation of the video game series, closely resembles a cormorant in both design and name. The cormorant was chosen as the emblem for the Ministry of Defence Joint Services Command and Staff College at Shrivenham. A bird famed for flight, sea fishing and land nesting was felt to be particularly appropriate for a college that unified leadership training and development for the Army, Navy and Royal Air Force. After a member produced a mock magazine cover from a photograph of roosting cormorants, the bird became the unofficial mascot of the Pentax Discuss Mailing List with many posts dedicated to discussion of the photography of the species.	Biology and health sciences	Pelecanimorphae	null

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