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210892 | https://en.wikipedia.org/wiki/Firefox | Firefox | Mozilla Firefox, or simply Firefox, is a free and open source web browser developed by the Mozilla Foundation and its subsidiary, the Mozilla Corporation. It uses the Gecko rendering engine to display web pages, which implements current and anticipated web standards. Firefox is available for Windows 10 and later versions of Windows, macOS, and Linux. Its unofficial ports are available for various Unix and Unix-like operating systems, including FreeBSD, OpenBSD, NetBSD, and other operating systems, such as reactOS. Firefox is also available for Android and iOS. However, as with all other iOS web browsers, the iOS version uses the WebKit layout engine instead of Gecko due to platform requirements. An optimized version is also available on the Amazon Fire TV as one of the two main browsers available with Amazon's Silk Browser.
Firefox is the spiritual successor of Netscape Navigator, as the Mozilla community was created by Netscape in 1998, before its acquisition by AOL. Firefox was created in 2002 under the codename "Phoenix" by members of the Mozilla community who desired a standalone browser rather than the Mozilla Application Suite bundle. During its beta phase, it proved to be popular with its testers and was praised for its speed, security, and add-ons compared to Microsoft's then-dominant Internet Explorer 6. It was released on November 9, 2004, and challenged Internet Explorer's dominance with 60 million downloads within nine months. In November 2017, Firefox began incorporating new technology under the code name "Quantum" to promote parallelism and a more intuitive user interface.
Firefox usage share grew to a peak of 32.21% in November 2009, with Firefox 3.5 overtaking Internet Explorer 7, although not all versions of Internet Explorer as a whole; its usage then declined in competition with Google Chrome. , according to StatCounter, it had a 6.5% usage share on traditional PCs (i.e. as a desktop browser), making it the fourth-most popular PC web browser after Google Chrome (65%), Microsoft Edge (13%), and Safari (9.0%).
History
The project began as an experimental branch of the Mozilla project by Dave Hyatt, Joe Hewitt, and Blake Ross. They believed the commercial requirements of Netscape's sponsorship and developer-driven feature creep compromised the utility of the Mozilla browser. To combat what they saw as the Mozilla Suite's software bloat, they created a standalone browser, with which they intended to replace the Mozilla Suite. Version 0.1 was released on September 23, 2002. On April 3, 2003, the Mozilla Organization announced that it planned to change its focus from the Mozilla Suite to Firefox and Thunderbird.
The Firefox project has undergone several name changes. The nascent browser was originally named Phoenix, after the mythical bird that rose triumphantly from the ashes of its dead predecessor (in this case, from the "ashes" of Netscape Navigator, after it was sidelined by Microsoft Internet Explorer in the "First Browser War"). Phoenix was renamed in 2003 due to a trademark claim from Phoenix Technologies. The replacement name, Firebird, provoked an intense response from the Firebird database software project. The Mozilla Foundation reassured them that the browser would always bear the name Mozilla Firebird to avoid confusion. After further pressure, Mozilla Firebird became Mozilla Firefox on February 9, 2004. The name Firefox was said to be derived from a nickname of the red panda, which became the mascot for the newly named project. For the abbreviation of Firefox, Mozilla prefers Fx or fx, although it is often abbreviated as FF or Ff.
The Firefox project went through many versions before version 1.0 and had already gained a great deal of acclaim from numerous media outlets, such as Forbes and The Wall Street Journal. Among Firefox's popular features were the integrated pop-up blocker, tabbed browsing, and an extension mechanism for adding functionality. Although these features have already been available for some time in other browsers such as the Mozilla Suite and Opera, Firefox was the first of these browsers to have achieved large-scale adoption so quickly. Firefox attracted attention as an alternative to Internet Explorer, which had come under fire for its alleged poor program design and insecurity—detractors cite IE's lack of support for certain Web standards, use of the potentially dangerous ActiveX component, and vulnerability to spyware and malware installation. Microsoft responded by releasing Windows XP Service Pack 2, which added several important security features to Internet Explorer 6.
Version 1.0 of Firefox was released on November 9, 2004. This was followed by version 1.5 in November 2005, version 2.0 in October 2006, version 3.0 in June 2008, version 3.5 in June 2009, version 3.6 in January 2010, and version 4.0 in March 2011. From version 5 onwards, the development and release model changed into a "rapid" one; by the end of 2011 the stable release was version 9, and by the end of 2012 it reached version 17.
Major redesigns of its graphical user interface occurred on versions 4.0 (Strata) in March 2011, 29.0 (Australis) in April 2014, 57.0 (Photon) in November 2017, and 89.0 (Proton) in June 2021.
In 2016, Mozilla announced a project known as Quantum, which sought to improve Firefox's Gecko engine and other components to improve the browser's performance, modernize its architecture, and transition the browser to a multi-process model. These improvements came in the wake of decreasing market share to Google Chrome, as well as concerns that its performance was lapsing in comparison. Despite its improvements, these changes required existing add-ons for Firefox to be made incompatible with newer versions, in favor of a new extension system that is designed to be similar to Chrome and other recent browsers. Firefox 57, which was released in November 2017, was the first version to contain enhancements from Quantum, and has thus been named Firefox Quantum. A Mozilla executive stated that Quantum was the "biggest update" to the browser since version 1.0. Unresponsive and crashing pages only affect other pages loaded within the same process. While Chrome uses separate processes for each loaded tab, Firefox distributes tabs over four processes by default (since Quantum), in order to balance memory consumption and performance. The process count can be adjusted, where more processes increase performance at the cost of memory, therefore suitable for computers with larger RAM capacity.
On May 3, 2019, the expiry of an intermediate signing certificate on Mozilla servers caused Firefox to automatically disable and lock all browser extensions (add-ons). Mozilla began the roll-out of a fix shortly thereafter, using their Mozilla Studies component.
Support for Adobe Flash was dropped on January 6, 2021, with the release of Firefox 85.
On June 1, 2021, Firefox's 'Proton' redesign was offered through its stable release channel after being made available in the beta builds. While users were initially allowed to revert to the old design through about:config, the corresponding key-value pairs reportedly stopped working in later builds, resulting in criticism. These included accessibility concerns despite Mozilla's claim to "continue to work with the accessibility community" and had not been resolved .
On January 13, 2022, an issue with Firefox's HTTP/3 implementation resulted in a widespread outage for several hours.
On September 26, 2023, Firefox 118.0 introduced on-device translation of web page content.
On January 23, 2024, along with the release of Firefox 122.0, Mozilla introduced an official APT repository for Debian-based Linux distributions.
Features
Features of the desktop edition include tabbed browsing, full-screen mode, spell checking, incremental search, smart bookmarks, bookmarking and downloading through drag and drop, a download manager, user profile management, private browsing, bookmark tags, bookmark exporting, offline mode, a screenshot tool, web development tools, a "page info" feature which shows a list of page metadata and multimedia items, a configuration menu at about:config for power users, and location-aware browsing (also known as "geolocation") based on a Google service. Firefox has an integrated search system which uses Google by default in most markets. DNS over HTTPS is another feature whose default behaviour is determined geographically.
Firefox provides an environment for web developers in which they can use built-in tools, such as the Error Console or the DOM Inspector, and extensions, such as Firebug and more recently there has been an integration feature with Pocket. Firefox Hello was an implementation of WebRTC, added in October 2014, which allows users of Firefox and other compatible systems to have a video call, with the extra feature of screen and file sharing by sending a link to each other. Firefox Hello was scheduled to be removed in September 2016.
Former features include a File Transfer Protocol (FTP) client for browsing file servers, the ability to block images from individual domains (until version 72), a 3D page inspector (versions 11 to 46), tab grouping (until version 44), and the ability to add customized extra toolbars (until version 28).
Browser extensions
Functions can be added through add-ons created by third-party developers. Add-ons are primarily coded using an HTML, CSS, JavaScript, with API known as WebExtensions, which is designed to be compatible with Google Chrome and Microsoft Edge extension systems. Firefox previously supported add-ons using the XUL and XPCOM APIs, which allowed them to directly access and manipulate much of the browser's internal functionality. As compatibility was not included in the multi-process architecture, XUL add-ons have been deemed Legacy add-ons and are no longer supported on Firefox 57 "Quantum" and newer.
Mozilla has occasionally installed extensions for users without their permission. This happened in 2017 when an extension designed to promote the show Mr. Robot was silently added in an update to Firefox.
Themes
Firefox can have themes added to it, which users can create or download from third parties to change the appearance of the browser. Firefox also provides dark, light, and system themes.
Guest session
In 2013, Firefox for Android added a guest session mode, which wiped browsing data such as tabs, cookies, and history at the end of each guest session. Guest session data was kept even when restarting the browser or device, and deleted only upon a manual exit. The feature was removed in 2019, purportedly to "streamline the experience".
Standards
Firefox implements many web standards, including HTML4 (almost full HTML5), XML, XHTML, MathML, SVG 1.1 (full), SVG 2 (partial), CSS (with extensions), ECMAScript (JavaScript), DOM, XSLT, XPath, and APNG (Animated PNG) images with alpha transparency. Firefox also implements standards proposals created by the WHATWG such as client-side storage, and the canvas element. These standards are implemented through the Gecko layout engine, and SpiderMonkey JavaScript engine. Firefox 4 was the first release to introduce significant HTML5 and CSS3 support.
Firefox has passed the Acid2 standards-compliance test since version 3.0. Mozilla had originally stated that they did not intend for Firefox to pass the Acid3 test fully because they believed that the SVG fonts part of the test had become outdated and irrelevant, due to WOFF being agreed upon as a standard by all major browser makers. Because the SVG font tests were removed from the Acid3 test in September 2011, Firefox 4 and greater scored 100/100.
Firefox also implements "Safe Browsing," a proprietary protocol from Google used to exchange data related with phishing and malware protection.
Firefox supports the playback of video content protected by HTML5 Encrypted Media Extensions (EME), since version 38. For security and privacy reasons, EME is implemented within a wrapper of open-source code that allows execution of a proprietary DRM module by Adobe Systems—Adobe Primetime Content Decryption Module (CDM). CDM runs within a "sandbox" environment to limit its access to the system and provide it a randomized device ID to prevent services from uniquely identifying the device for tracking purposes. The DRM module, once it has been downloaded, is enabled, and disabled in the same manner as other plug-ins. Since version 47, "Google's Widevine CDM on Windows and Mac OS X so streaming services like Amazon Video can switch from Silverlight to encrypted HTML5 video" is also supported. Mozilla justified its partnership with Adobe and Google by stating:
and that it is "an important step on Mozilla's roadmap to remove NPAPI plugin support." Upon the introduction of EME support, builds of Firefox on Windows were also introduced that exclude support for EME. The Free Software Foundation and Cory Doctorow condemned Mozilla's decision to support EME.
Security
From its inception, Firefox was positioned as a security-focused browser. At the time, Internet Explorer, the dominant browser, was facing a security crisis. Multiple vulnerabilities had been found, and malware like Download.Ject could be installed simply by visiting a compromised website. The situation was so bad that the US Government issued a warning against using Internet Explorer. Firefox, being less integrated with the operating system, was considered a safer alternative since it was less likely to have issues that could completely compromise a computer. This led to a significant increase in Firefox's popularity during the early 2000s as a more secure alternative. Moreover, Firefox was considered to have fewer actively exploitable security vulnerabilities compared to its competitors. In 2006, The Washington Post reported that exploit code for known security vulnerabilities in Internet Explorer were available for 284 days compared to only nine days for Firefox before the problem was fixed. A Symantec study around the same period showed that even though Firefox had a higher number of vulnerabilities, on average vulnerabilities were fixed faster in Firefox than in other browsers during that period.
During this period, Firefox used a monolithic architecture, like most browsers at the time. This meant all browser components ran in a single process with access to all system resources. This setup had multiple security issues. If a web page used too many resources, the entire Firefox process would hang or crash, affecting all tabs. Additionally, any exploit could easily access system resources, including user files. Between 2008 and 2012, most browsers shifted to a multiprocess architecture, isolating high-risk processes like rendering, media, GPU, and networking. However, Firefox was slower to adopt this change. It wasn't until 2015 that Firefox started its Electrolysis (e10s) project to implement sandboxing across multiple components. This rewrite relied on interprocess communication using Chromium's interprocess communication library and placed various component including the rendering component in its own sandbox. Firefox released this rewrite in to beta in August 2016, noting a 10–20% increase in memory usage, which was lower than Chrome's at the time. However, the rewrite caused issues with their legacy extension API, which was not designed to work cross-process and required shim code to function correctly. After over a year in beta, the rewrite was enabled by default all users of Firefox in November 2017.
In 2012, Mozilla launched a new project called Servo to write a completely new and experimental browser engine utilizing memory safe techniques written in Rust. In 2018, Mozilla opted to integrate parts of the Servo project into the Gecko engine in a project codenamed the Quantum project. The project completely overhauled Firefox's page rendering code resulting in performance and stability gains while also improving the security of existing components. Additionally, the older incompatible extension API was removed in favour of a WebExtension API that more closely resembled Google Chrome's extension system. This broke compatibility with older extensions but resulted in lesser vulnerabilities and a much more maintainable extension system. While the Servo project was intended to replace more parts of the Gecko Engine, this plan never came to fruition. In 2020, Mozilla laid off all developers on the Servo team transferring ownership of the project to the Linux Foundation.
Privacy
When Firefox initially released, it used a custom script permission policy where scripts that were signed by the page could gain access to higher privilege actions such as the ability to set a user's preferences. However, this model was not widely used and was later discontinued by Firefox. Modern day Firefox instead follows the standard same-origin policy permission model that is followed by most modern browsers which disallows scripts from accessing any privileged data including data about other websites.
It uses TLS to protect communications with web servers using strong cryptography when using the HTTPS protocol. The freely available HTTPS Everywhere add-on enforces HTTPS, even if a regular HTTP URL is entered. Firefox now supports HTTP/2.
In February 2013, plans were announced for Firefox 22 to disable third-party cookies by default. However, the introduction of the feature was then delayed so Mozilla developers could "collect and analyze data on the effect of blocking some third-party cookies." Mozilla also collaborated with Stanford University's "Cookie Clearinghouse" project to develop a blacklist and whitelist of sites that will be used in the filter.
Version 23, released in August 2013, followed the lead of its competitors by blocking iframe, stylesheet, and script resources served from non-HTTPS servers embedded on HTTPS pages by default. Additionally, JavaScript could also no longer be disabled through Firefox's preferences, and JavaScript was automatically re-enabled for users who upgraded to 23 or higher with it disabled. The change was made due to the fact the JavaScript was being used across a majority of websites on the web and disabling JavaScript could potentially have untoward repercussions on inexperienced users who are unaware of its impact. Firefox also cited the fact that extensions like NoScript, that can disable JavaScript in a more controlled fashion, were widely available. The following release added the ability to disable JavaScript through the developer tools for testing purposes.
Beginning with Firefox 48, all extensions must be signed by Mozilla to be used in release and beta versions of Firefox. Firefox 43 blocked unsigned extensions but allowed enforcement of extension signing to be disabled. All extensions must be submitted to Mozilla Add-ons and be subject to code analysis in order to be signed, although extensions do not have to be listed on the service to be signed. On May 2, 2019, Mozilla announced that it would be strengthening the signature enforcement with methods that included the retroactive disabling of old extensions now deemed to be insecure.
Since version 60 Firefox includes the option to use DNS over HTTPS (DoH), which causes DNS lookup requests to be sent encrypted over the HTTPS protocol. To use this feature the user must set certain preferences beginning with "network.trr" (Trusted Recursive Resolver) in about:config: if network.trr.mode is 0, DoH is disabled; 1 activates DoH in addition to unencrypted DNS; 2 causes DoH to be used before unencrypted DNS; to use only DoH, the value must be 3. By setting network.trr.uri to the URL, special Cloudflare servers will be activated. Mozilla has a privacy agreement with this server host that restricts their collection of information about incoming DNS requests.
On May 21, 2019, Firefox was updated to include the ability to block scripts that used a computer's CPU to mine cryptocurrency without a user's permission, in Firefox version 67.0. The update also allowed users to block known fingerprinting scripts that track their activity across the web, however it does not resist fingerprinting on its own.
In March 2021, Firefox launched SmartBlock in version 87 to offer protection against cross-site tracking, without breaking the websites users visit. Also known as state partitioning or "total cookie protection", SmartBlock works via a feature in the browser that isolates data from each site visited by the user to ensure that cross-site scripting is very difficult if not impossible. The feature also isolates local storage, service workers and other common ways for sites to store data.
Localizations
Firefox is a widely localized web browser. Mozilla uses the in-house Pontoon localization platform. The first official release in November 2004 was available in 24 different languages and for 28 locales. In 2019, Mozilla released Project Fluent a localization system that allows translators to be more flexible with their translation than to be constrained in one-to-one translation of strings. the supported versions of Firefox are available in 97 locales (88 languages).
Platform availability
There are desktop versions of Firefox for Microsoft Windows, macOS, and Linux, while Firefox for Android is available for Android (formerly Firefox for mobile, it also ran on Maemo, MeeGo and Firefox OS) and Firefox for iOS is available for iOS. Smartphones that support Linux but not Android, or iOS apps can also run Firefox in its desktop version, for example using postmarketOS.
Firefox source code may be compiled for various operating systems; however, officially distributed binaries are provided for the following:
Microsoft Windows
Firefox 1.0 was released for Windows 95, as well as Windows NT 4.0 or later. Some users reported the 1.x builds were operable (but not installable) on Windows NT 3.51.
The version 42.0 release includes the first x64 build. It required Windows 7 or Server 2008 R2. Starting from version 49.0, Firefox for Windows requires and uses the SSE2 instruction set.
In September 2013, Mozilla released a Metro-style version of Firefox, optimized for touchscreen use, on the "Aurora" release channel. However, on March 14, 2014, Mozilla cancelled the project because of a lack of user adoption.
In March 2017, Firefox 52 ESR, the last version of the browser for Windows XP and Windows Vista, was released. Support for Firefox 52 ESR ended in June 2018.
Traditionally, installing the Windows version of Firefox entails visiting the Firefox website and downloading an installer package, depending on the desired localization and system architecture. In November 2021, Mozilla made Firefox available on Microsoft Store. The Store-distributed package does not interfere with the traditional installation.
The last version of Firefox for Windows 7 and 8 is Firefox 115 ESR, which was released in July 2023. Its end-of-life was initially planned to be in October 2024, however in July 2024, a Mozilla employee announced in a comment on Reddit that the company consider extending the support beyond the initial date, the duration of that extension being yet to be defined. In September 2024, the extension was announced for an initial period of six months, leading the end-of-life date to March 2025. In another Reddit comment, another Mozilla employee said that despite the extension "isn't going to be free [for Mozilla]", there was enough users to justify extending the support. In the release calendar page, a note states that Mozilla will re-evaluate the situation in early 2025 to see if another extension will be needed or not and statute about 115 ESR end-of-life then.
macOS
The first official release (Firefox version 1.0) supported macOS (then called Mac OS X) on the PowerPC architecture. Mac OS X builds for the IA-32 architecture became available via a universal binary which debuted with Firefox 1.5.0.2 in 2006.
Starting with version 4.0, Firefox was released for the x64 architecture to which macOS had migrated. Version 4.0 also dropped support for PowerPC architecture, although other projects continued development of a PowerPC version of Firefox.
Firefox was originally released for Mac OS X 10.0 and higher. The minimum OS then increased to Mac OS X 10.2 in Firefox 1.5 and 10.4 in Firefox 3. Firefox 4 dropped support for Mac OS X 10.4 and PowerPC Macs, and Firefox 17 dropped support for Mac OS X 10.5 entirely. The system requirements were left unchanged until 2016, when Firefox 49 dropped support for Mac OS X 10.6–10.8. Mozilla ended support for OS X 10.9–10.11 in Firefox 79, with those users being supported on the Firefox 78 ESR branch until November 2021. Most recently, Mozilla ended support for macOS 10.12–10.14 in Firefox 116, with those users being supported on the Firefox 115 ESR branch until late 2024. In September 2024 however, an extension was announced for the 115 ESR branch for an initial period of six months, leading the end-of-life date to March 2025.
Linux
Since its inception, Firefox for Linux supported the 32-bit memory architecture of the IA-32 instruction set. 64-bit builds were introduced in the 4.0 release. The 46.0 release replaced GTK 2.18 with 3.4 as a system requirement on Linux and other systems running X.Org. Starting with 53.0, the 32-bit builds require the SSE2 instruction set.
Firefox for Android
Firefox for mobile, code-named "Fennec", was first released for Maemo in January 2010 with version 1.0 and for Android in March 2011 with version 4.0. Support for Maemo was discontinued after version 7, released in September 2011. Fennec had a user interface optimized for phones and tablets. It included the Awesome Bar, tabbed browsing, add-on support, a password manager, location-aware browsing, and the ability to synchronize with the user's other devices with Mozilla Firefox using Firefox Sync. At the end of its existence, it had a market share of 0.5% on Android.
In August 2020, Mozilla launched a new version of its Firefox for Android app, named Firefox Daylight to the public and codenamed Fenix, after a little over a year of testing. It boasted higher speeds with its new GeckoView engine, which is described as being "the only independent web engine browser available on Android". It also added Enhanced Tracking Protection 2.0, a feature that blocks many known trackers on the Internet. It also added the ability to place the address bar on the bottom, and a new Collections feature. However, it was criticized for only having nine Add-ons at launch, and missing certain features. In response, Mozilla stated that they will allow more Add-ons with time.
Firefox for iOS
Mozilla initially refused to port Firefox to iOS, due to the restrictions Apple imposed on third-party iOS browsers. Instead of releasing a full version of the Firefox browser, Mozilla released Firefox Home, a companion app for the iPhone and iPod Touch based on the Firefox Sync technology, which allowed users to access their Firefox browsing history, bookmarks, and recent tabs. It also included Firefox's "Awesomebar" location bar. Firefox Home was not a web browser, the application launched web pages in either an embedded viewer for that one page, or by opening the page in the Safari app. Mozilla pulled Firefox Home from the App Store in September 2012, stating it would focus its resources on other projects. The company subsequently released the source code of Firefox Home's underlying synchronization software.
In April 2013, then-Mozilla CEO Gary Kovacs said that Firefox would not come to iOS if Apple required the use of the WebKit layout engine to do so. One reason given by Mozilla was that prior to iOS 8, Apple had supplied third-party browsers with an inferior version of their JavaScript engine which hobbled their performance, making it impossible to match Safari's JavaScript performance on the iOS platform. Apple later opened their "Nitro" JavaScript engine to third-party browsers. In 2015, Mozilla announced it was moving forward with Firefox for iOS, with a preview release made available in New Zealand in September of that year. It was fully released in November later that year. It is the first Firefox-branded browser not to use the Gecko layout engine as is used in Firefox for desktop and mobile. Apple's policies require all iOS apps that browse the web to use the built-in WebKit rendering framework and WebKit JavaScript, so using Gecko is not possible. Unlike Firefox on Android, Firefox for iOS does not support browser add-ons.
In November 2016, Firefox released a new iOS app titled Firefox Focus, a private web browser.
Firefox Reality (AR/VR)
Firefox Reality was released for Augmented Reality and Virtual Reality headsets in September 2018. It supports traditional web-browsing through 2D windows and immersive VR pages through Web VR. Firefox Reality is available on HTC Vive, Oculus, Google Daydream and Microsoft Hololens headsets. In February 2022 Mozilla announced that Igalia took over stewardship of this project under the new name of Wolvic.
Unofficial ports
Firefox has also been ported to FreeBSD, NetBSD, OpenBSD, OpenIndiana, OS/2, ArcaOS, SkyOS, RISC OS and BeOS/Haiku, and an unofficial rebranded version called Timberwolf has been available for AmigaOS 4. An unofficial continuation of the Mac OS X PowerPC release was actively developed as TenFourFox until October 5, 2021.
The Firefox port for OpenBSD is maintained by Landry Breuil since 2010. Firefox is regularly built for the current branch of the operating system, the latest versions are packaged for each release and remain frozen until the next release. In 2017, Landry began hosting packages of newer Firefox versions for OpenBSD releases from 6.0 onwards, making them available to installations without the ports system.
The Solaris port of Firefox (including OpenSolaris) was maintained by the Oracle Solaris Desktop Beijing Team, until March 2017 when the team was disbanded. There was also an unofficial port of Firefox 3.6.x to IBM AIX and of v1.7.x to UnixWare.
Channels and release schedule
In March 2011, Mozilla presented plans to switch to the rapid release model, a faster 16-week development cycle, similar to Google Chrome. Ars Technica noted that this new cycle entailed "significant technical and operational challenges" for Mozilla (notably preserving third-party add-on compatibility), but that it would help accelerate Firefox's adoption of new web standards, feature, and performance improvements. This plan was implemented in April 2011. The release process was split into four "channels", with major releases trickling down to the next channel every six to eight weeks. For example, the Nightly channel would feature a preliminary unstable version of Firefox 6, which would move to the experimental "Aurora" channel after preliminary testing, then to the more stable "beta" channel, before finally reaching the public release channel, with each stage taking around six weeks. For corporations, Mozilla introduced an Extended Support Release (ESR) channel, with new versions released every 30 weeks (and supported for 12 more weeks after a new ESR version is released), though Mozilla warned that it would be less secure than the release channel, since security patches would only be backported for high-impact vulnerabilities.
In 2017, Mozilla abandoned the Aurora channel, which saw low uptake, and rebased Firefox Developer Edition onto the beta channel. Mozilla uses A/B testing and a staged rollout mechanism for the release channel, where updates are first presented to a small fraction of users, with Mozilla monitoring its telemetry for increased crashes or other issues before the update is made available to all users. In 2020, Firefox moved to a four-week release cycle, to catch up with Chrome in support for new web features. Chrome switched to a four-week cycle a year later.
Licensing
Firefox source code is free software, with most of it being released under the Mozilla Public License (MPL) version 2.0. This license permits anyone to view, modify, or redistribute the source code. As a result, several publicly released applications have been built from it, including Firefox's predecessor Netscape, the customizable Pale Moon, and the privacy focused Tor Browser.
In the past, Firefox was licensed solely under the MPL, then version 1.1, which the Free Software Foundation criticized for being weak copyleft, as the license permitted, in limited ways, proprietary derivative works. Additionally, code only licensed under MPL 1.1 could not legally be linked with code under the GPL. To address these concerns, Mozilla re-licensed most of Firefox under the tri-license scheme of MPL 1.1, GPL 2.0, or LGPL 2.1. Since the re-licensing, developers were free to choose the license under which they received most of the code, to suit their intended use: GPL or LGPL linking and derivative works when one of those licenses is chosen, or MPL use (including the possibility of proprietary derivative works) if they chose the MPL. However, on January 3, 2012, Mozilla released the GPL-compatible MPL 2.0, and with the release of Firefox 13 on June 5, 2012, Mozilla used it to replace the tri-licensing scheme.
Trademark and logo
The name "Mozilla Firefox" is a registered trademark of Mozilla; along with the official Firefox logo, it may only be used under certain terms and conditions. Anyone may redistribute the official binaries in unmodified form and use the Firefox name and branding for such distribution, but restrictions are placed on distributions which modify the underlying source code. The name "Firefox" derives from a nickname of the red panda. Mozilla celebrated Red Pandas.
Mozilla has placed the Firefox logo files under open-source licenses, but its trademark guidelines do not allow displaying altered or similar logos in contexts where trademark law applies.
There has been some controversy over the Mozilla Foundation's intentions in stopping certain open-source distributions from using the "Firefox" trademark. Open-source browsers "enable greater choice and innovation in the market rather than aiming for mass-market domination." Mozilla Foundation Chairperson Mitchell Baker explained in an interview in 2007 that distributions could freely use the Firefox trademark if they did not modify source code, and that the Mozilla Foundation's only concern was with users getting a consistent experience when they used "Firefox".
To allow distributions of the code without using the official branding, the Firefox build system contains a "branding switch". This switch, often used for alphas ("Auroras") of future Firefox versions, allows the code to be compiled without the official logo and name and can allow a derivative work unencumbered by restrictions on the Firefox trademark to be produced. In the unbranded build, the trademarked logo and name are replaced with a freely distributable generic globe logo and the name of the release series from which the modified version was derived.
Distributing modified versions of Firefox under the "Firefox" name required explicit approval from Mozilla for the changes made to the underlying code, and required the use of all of the official branding. For example, it was not permissible to use the name "Firefox" without also using the official logo. When the Debian project decided to stop using the official Firefox logo in 2006 (because Mozilla's copyright restrictions at the time were incompatible with Debian's guidelines), they were told by a representative of the Mozilla Foundation that this was not acceptable and was asked either to comply with the published trademark guidelines or cease using the "Firefox" name in their distribution. Debian switched to branding their modified version of Firefox "Iceweasel" (but in 2016 switched back to Firefox), along with other Mozilla software. GNU IceCat is another derived version of Firefox distributed by the GNU Project, which maintains its separate branding.
Branding and visual identity
The Firefox icon is a trademark used to designate the official Mozilla build of the Firefox software and builds of official distribution partners. For this reason, software distributors who distribute modified versions of Firefox do not use the icon.
Early Firebird and Phoenix releases of Firefox were considered to have reasonable visual designs but fell short when compared to many other professional software packages. In October 2003, professional interface designer Steven Garrity authored an article covering everything he considered to be wrong with Mozilla's visual identity.
Shortly afterwards, the Mozilla Foundation invited Garrity to head up the new visual identity team. The release of Firefox 0.8 in February 2004 saw the introduction of the new branding efforts. Included were new icon designs by silverorange, a group of web developers with a long-standing relationship with Mozilla. The final renderings are by Jon Hicks, who had worked on Camino. The logo was later revised and updated, fixing several flaws found when it was enlarged. The animal shown in the logo is a stylized fox, although "firefox" is usually a common name for the red panda. The panda, according to Hicks, "didn't really conjure up the right imagery" and was not widely known.
In June 2019, Mozilla unveiled a revised Firefox logo, which was officially implemented on version 70. The new logo is part of an effort to build a brand system around Firefox and its complementary apps and services, which are now being promoted as a suite under the Firefox brand.
Promotion
Firefox was adopted rapidly, with 100 million downloads in its first year of availability. This was followed by a series of aggressive marketing campaigns starting in 2004 with a series of events Blake Ross and Asa Dotzler called "marketing weeks".
Firefox continued to heavily market itself by releasing a marketing portal dubbed "Spread Firefox" (SFX) on September 12, 2004. It debuted along with the Firefox Preview Release, creating a centralized space for the discussion of various marketing techniques. The release of their manifesto stated that "the Mozilla project is a global community of people who believe that openness, innovation and opportunity are key to the continued health of the Internet." A two-page ad in the edition of December 16 of The New York Times, placed by Mozilla Foundation in coordination with Spread Firefox, featured the names of the thousands of people worldwide who contributed to the Mozilla Foundation's fundraising campaign to support the launch of the Firefox 1.0 web browser. SFX portal enhanced the "Get Firefox" button program, giving users "referrer points" as an incentive. The site lists the top 250 referrers. From time to time, the SFX team or SFX members launch marketing events organized at the Spread Firefox website. As a part of the Spread Firefox campaign, there was an attempt to break the world download record with the release of Firefox 3. This resulted in an official certified Guinness world record, with over eight million downloads. In February 2011, Mozilla announced that it would be retiring Spread Firefox (SFX). Three months later, in May 2011, Mozilla officially closed Spread Firefox. Mozilla wrote that "there are currently plans to create a new iteration of this website [Spread Firefox] at a later date."
In celebration of the third anniversary of the founding of the Mozilla Foundation, the "World Firefox Day" campaign was established on July 15, 2006, and ran until September 15, 2006. Participants registered themselves and a friend on the website for nomination to have their names displayed on the Firefox Friends Wall, a digital wall that was displayed at the headquarters of the Mozilla Foundation.
The Firefox community has also engaged in the promotion of their web browser. In 2006, some of Firefox's contributors from Oregon State University made a crop circle of the Firefox logo in an oat field near Amity, Oregon, near the intersection of Lafayette Highway and Walnut Hill Road. After Firefox reached 500 million downloads on February 21, 2008, the Firefox community celebrated by visiting Freerice to earn 500 million grains of rice.
Other initiatives included Live Chat – a service Mozilla launched in 2007 that allowed users to seek technical support from volunteers. The service was later retired.
To promote the launch of Firefox Quantum in November 2017, Mozilla partnered with Reggie Watts to produce a series of TV ads and social media content.
Performance
2000s
In December 2005, Internet Week ran an article in which many readers reported high memory usage in Firefox 1.5. Mozilla developers said that the higher memory use of Firefox 1.5 was at least partially due to the new fast backwards-and-forwards (FastBack) feature. Other known causes of memory problems were malfunctioning extensions such as Google Toolbar and some older versions of AdBlock, or plug-ins, such as older versions of Adobe Acrobat Reader. When PC Magazine in 2006 compared memory usage of Firefox 2, Opera 9, and Internet Explorer 7, they found that Firefox used approximately as much memory as each of the other two browsers.
In 2006, Softpedia noted that Firefox 1.5 took longer to start up than other browsers, which was confirmed by further speed tests.
Internet Explorer 6 launched more swiftly than Firefox 1.5 on Windows XP since many of its components were built into the OS and loaded during system startup. As a workaround for the issue, a preloader application was created that loaded components of Firefox on startup, similar to Internet Explorer. A Windows Vista feature called SuperFetch performs a similar task of preloading Firefox if it is used often enough.
Tests performed by PC World and Zimbra in 2006 indicated that Firefox 2 used less memory than Internet Explorer 7. Firefox 3 used less memory than Internet Explorer 7, Opera 9.50 Beta, Safari 3.1 Beta, and Firefox 2 in tests performed by Mozilla, CyberNet, and The Browser World. In mid-2009, BetaNews benchmarked Firefox 3.5 and declared that it performed "nearly ten times better on XP than Microsoft Internet Explorer 7".
2010s
In January 2010, Lifehacker compared the performance of Firefox 3.5, Firefox 3.6, Google Chrome 4 (stable and Dev versions), Safari 4, and Opera (10.1 stable and 10.5 pre-alpha versions). Lifehacker timed how long browsers took to start and reach a page (both right after boot-up and after running at least once already), timed how long browsers took to load nine tabs at once, tested JavaScript speeds using Mozilla's Dromaeo online suite (which implements Apple's SunSpider and Google's V8 tests) and measured memory usage using Windows 7's process manager. They concluded that Firefox 3.5 and 3.6 were the fifth- and sixth-fastest browsers, respectively, on startup, 3.5 was third- and 3.6 was sixth-fastest to load nine tabs at once, 3.5 was sixth- and 3.6 was fifth-fastest on the JavaScript tests. They also concluded that Firefox 3.6 was the most efficient with memory usage followed by Firefox 3.5.
In February 2012, Tom's Hardware performance tested Chrome 17, Firefox 10, Internet Explorer 9, Opera 11.61, and Safari 5.1.2 on Windows 7. Tom's Hardware summarized their tests into four categories: Performance, Efficiency, Reliability, and Conformance. In the performance category they tested HTML5, Java, JavaScript, DOM, CSS 3, Flash, Silverlight, and WebGL (WebGL 2 is current as of version 51; and Java and Silverlight stop working as of version 52)—they also tested startup time and page load time. The performance tests showed that Firefox was either "acceptable" or "strong" in most categories, winning three categories (HTML5, HTML5 hardware acceleration, and Java) only finishing "weak" in CSS performance. In the efficiency tests, Tom's Hardware tested memory usage and management. With this category, it determined that Firefox was only "acceptable" at performing light memory usage, while it was "strong" at performing heavy memory usage. In the reliability category, Firefox performed a "strong" amount of proper page loads. For the final category, conformance, it was determined that Firefox had "strong" conformance for JavaScript and HTML5. So in conclusion, Tom's Hardware determined that Firefox was the best browser for Windows 7 OS, but that it only narrowly beat Google Chrome.
In June 2013, Tom's Hardware again performance tested Firefox 22, Chrome 27, Opera 12, and Internet Explorer 10. They found that Firefox slightly edged out the other browsers in their "performance" index, which examined wait times, JavaScript execution speed, HTML5/CSS3 rendering, and hardware acceleration performance. Firefox also scored the highest on the "non-performance" index, which measured memory efficiency, reliability, security, and standards conformance, finishing ahead of Chrome, the runner-up. Tom's Hardware concluded by declaring Firefox the "sound" winner of the performance benchmarks.
In January 2014, a benchmark testing the memory usage of Firefox 29, Google Chrome 34, and Internet Explorer 11 indicated that Firefox used the least memory when a substantial number of tabs were open.
In benchmark testing in early 2015 on a "high-end" Windows machine, comparing Microsoft Edge [Legacy], Internet Explorer, Firefox, Chrome, and Opera, Firefox achieved the highest score on three of the seven tests. Four different JavaScript performance tests gave conflicting results. Firefox surpassed all other browsers on the Peacekeeper benchmark, but was behind the Microsoft products when tested with SunSpider. Measured with Mozilla's Kraken, it came second place to Chrome, while on Google's Octane challenge it took third behind Chrome and Opera. Firefox took the lead with WebXPRT, which runs several typical HTML5 and JavaScript tasks. Firefox, Chrome, and Opera all achieved the highest possible score on the Oort Online test, measuring WebGL rendering speed (WebGL 2 is now current). In terms of HTML5 compatibility testing, Firefox was ranked in the middle of the group.
A similar set of benchmark tests in 2016 showed Firefox's JavaScript performance on Kraken and the newer Jetstream tests trailing slightly behind all other tested browsers except Internet Explorer (IE), which performed relatively poorly. On Octane, Firefox came ahead of IE and Safari, but again slightly behind the rest, including Vivaldi and Microsoft Edge [Legacy]. Edge [Legacy] took overall first place on the Jetstream and Octane benchmarks.
Firefox Quantum
As of the adoption of Firefox 57 and Mozilla's Quantum project entering production browsers in November 2017, Firefox was tested to be faster than Chrome in independent JavaScript tests, and demonstrated to use less memory with many browser tabs opened. TechRadar rated it as the fastest web browser in a May 2019 report.
Usage share
Downloads have continued at an increasing rate since Firefox 1.0 was released, and Firefox had already been downloaded over one billion times. This number does not include downloads using software updates or those from third-party websites. They do not represent a user count, as one download may be installed on many machines, one person may download the software multiple times, or the software may be obtained from a third-party.
In July 2010, IBM asked all employees (about 400,000) to use Firefox as their default browser.
Firefox was the second-most used web browser until November 2011, when Google Chrome surpassed it. According to Mozilla, Firefox had more than 450 million users .
In October 2024, Firefox was the fourth-most widely used desktop browser, and it was the fourth-most popular with 2.95% of worldwide usage share of web browsers across all platforms.
According to the Firefox Public Data report by Mozilla, the active monthly count of Desktop clients has decreased from around 310 million in 2017 to 200 million in 2023.
From Oct 2020, the desktop market share of Firefox started to decline in countries where it used to be the most popular.
In Eritrea, it dropped from 50% in Oct 2020 to 9.32% in Sept 2021.
In Cuba, it dropped from 54.36% in Sept 2020 to 38.42% in Sept 2021.
The UK and US governments both follow the 2% rule. This states that only browsers with more than 2% market share among visitors of their websites will be supported. There are concerns that support for Firefox will be dropped because as of December 29, 2023, the browser market share among US government website visitors is 2.2%.
| Technology | Browsers | null |
210900 | https://en.wikipedia.org/wiki/Solar%20luminosity | Solar luminosity | The solar luminosity () is a unit of radiant flux (power emitted in the form of photons) conventionally used by astronomers to measure the luminosity of stars, galaxies and other celestial objects in terms of the output of the Sun.
One nominal solar luminosity is defined by the International Astronomical Union to be . The Sun is a weakly variable star, and its actual luminosity therefore fluctuates. The major fluctuation is the eleven-year solar cycle (sunspot cycle) that causes a quasi-periodic variation of about ±0.1%. Other variations over the last 200–300 years are thought to be much smaller than this.
Determination
Solar luminosity is related to solar irradiance (the solar constant). Solar irradiance is responsible for the orbital forcing that causes the Milankovitch cycles, which determine Earthly glacial cycles. The mean irradiance at the top of the Earth's atmosphere is sometimes known as the solar constant, . Irradiance is defined as power per unit area, so the solar luminosity (total power emitted by the Sun) is the irradiance received at the Earth (solar constant) multiplied by the area of the sphere whose radius is the mean distance between the Earth and the Sun:
where is the unit distance (the value of the astronomical unit in metres) and is a constant (whose value is very close to one) that reflects the fact that the mean distance from the Earth to the Sun is not exactly one astronomical unit.
| Physical sciences | Power | Basics and measurement |
211168 | https://en.wikipedia.org/wiki/Radio%20galaxy | Radio galaxy | A radio galaxy is a galaxy with giant regions of radio emission extending well beyond its visible structure. These energetic radio lobes are powered by jets from its active galactic nucleus. They have luminosities up to 1039 W at radio wavelengths between 10 MHz and 100 GHz. The radio emission is due to the synchrotron process. The observed structure in radio emission is determined by the interaction between twin jets and the external medium, modified by the effects of relativistic beaming. The host galaxies are almost exclusively large elliptical galaxies. Radio-loud active galaxies can be detected at large distances, making them valuable tools for observational cosmology. Recently, much work has been done on the effects of these objects on the intergalactic medium, particularly in galaxy groups and clusters.
The term "radio galaxy" is often used to refer to the entire jet system, rather than solely to its host galaxy. Some scientists consider the term "black hole jet system" more accurate and less confusing. Radio galaxies that reach the size of around 0.7 megaparsecs or more, are commonly called "giant radio galaxies".
Emission processes
The radio emission from radio-loud active galaxies is synchrotron emission, as inferred from its very smooth, broad-band nature and strong polarization. This implies that the radio-emitting plasma contains, at least, electrons with relativistic speeds (Lorentz factors of ~104) and magnetic fields. Since the plasma must be neutral, it must also contain either protons or positrons. There is no way of determining the particle content directly from observations of synchrotron radiation. Moreover, there is no way to determine the energy densities in particles and magnetic fields from observation: the same synchrotron emissivity may be a result of a few electrons and a strong field, or a weak field and many electrons, or something in between. It is possible to determine a minimum energy condition which is the minimum energy density that a region with a given emissivity can have, but for many years there was no particular reason to believe that the true energies were anywhere near the minimum energies.
A sister process to the synchrotron radiation is the inverse-Compton process, in which the relativistic electrons interact with ambient photons and Thomson scatter them to high energies. Inverse-Compton emission from radio-loud sources turns out to be particularly important in X-rays, and, because it depends only on the density of electrons, a detection of inverse-Compton scattering allows a somewhat model-dependent estimate of the energy densities in the particles and magnetic fields. This has been used to argue that many powerful sources are actually quite near the minimum-energy condition.
Synchrotron radiation is not confined to radio wavelengths: if the radio source can accelerate particles to high enough energies, features that are detected in the radio wavelengths may also be seen in the infrared, optical, ultraviolet or even X-ray. In the latter case the responsible electrons must have energies in excess of 1 TeV in typical magnetic field strengths. Again, polarization and continuum spectrum are used to distinguish the synchrotron radiation from other emission processes. Jets and hotspots are the usual sources of high-frequency synchrotron emission. It is hard to distinguish observationally between the synchrotron and inverse-Compton radiation, making them a subject of ongoing research.
Processes, collectively known as particle acceleration, produce populations of relativistic and non-thermal particles that give rise to synchrotron and inverse-Compton radiation. Fermi acceleration is one plausible particle acceleration process in radio-loud active galaxies.
Radio structures
Radio galaxies, and to a lesser extent, radio-loud quasars display a wide range of structures in radio maps. The most common large-scale structures are called lobes: these are double, often fairly symmetrical, roughly ellipsoidal structures placed on either side of the active nucleus. A significant minority of low-luminosity sources exhibit structures usually known as plumes which are much more elongated. Some radio galaxies show one or two long narrow features known as jets (the most famous example being the giant galaxy M87 in the Virgo cluster) coming directly from the nucleus and going to the lobes. Since the 1970s, the most widely accepted model has been that the lobes or plumes are powered by beams of high-energy particles and magnetic field coming from close to the active nucleus. The jets are believed to be the visible manifestations of the beams, and often the term jet is used to refer both to the observable feature and to the underlying flow.
In 1974, radio sources were divided by Bernard Fanaroff and Julia Riley into two classes, now known as Fanaroff and Riley Class I (FRI) and Class II (FRII). The distinction was originally made based on the morphology of the large-scale radio emission (the type was determined by the distance between the brightest points in the radio emission): FRI sources were brightest towards the centre, while FRII sources were brightest at the edges. Fanaroff and Riley observed that there was a reasonably sharp divide in luminosity between the two classes: FRIs were low-luminosity, FRIIs were high luminosity. With more detailed radio observations, the morphology turns out to reflect the method of energy transport in the radio source. FRI objects typically have bright jets in the centre, while FRIIs have faint jets but bright hotspots at the ends of the lobes. FRIIs appear to be able to transport energy efficiently to the ends of the lobes, while FRI beams are inefficient in the sense that they radiate a significant amount of their energy away as they travel.
In more detail, the FRI/FRII division depends on host-galaxy environment in the sense that the FRI/FRII transition appears at higher luminosities in more massive galaxies. FRI jets are known to be decelerating in the regions in which their radio emission is brightest, and so it seems that the FRI/FRII transition reflects whether a jet/beam can propagate through the host galaxy without being decelerated to sub-relativistic speeds by interaction with the intergalactic medium. From analysis of relativistic beaming effects, the jets of FRII sources are known to remain relativistic (with speeds of at least 0.5c) out to the ends of the lobes. The hotspots that are usually seen in FRII sources are interpreted as being the visible manifestations of shocks formed when the supersonic jet abruptly terminates at the end of the source, and their spectral energy distributions are consistent with this picture. Often multiple hotspots are seen, reflecting either continued outflow after the shock or movement of the jet termination point: the overall hotspot region is sometimes called the hotspot complex.
Names are given to several particular types of radio source based on their radio structure:
Classical double refers to an FRII source with clear hotspots.
Wide-angle tail normally refers to a source intermediate between standard FRI and FRII structure, with efficient jets and sometimes hotspots, but with plumes rather than lobes, found at or near the centres of clusters.
Narrow-angle tail or Head-tail source describes an FRI that appears to be bent by ram pressure as it moves through a cluster.
Fat doubles are sources with diffuse lobes but neither jets nor hotspots. Some such sources may be relics whose energy supply has been permanently or temporarily turned off.
Life cycles and dynamics
The largest radio galaxies have lobes or plumes extending to megaparsec scales (more in the case of giant radio galaxies like 3C236), implying a timescale for growth of the order of tens to hundreds of millions of years. This means that, except in the case of very small, very young sources, we cannot observe radio source dynamics directly, and so must resort to theory and inferences from large numbers of objects. Clearly, radio sources must start small and grow larger. In the case of sources with lobes, the dynamics are fairly simple: the jets feed the lobes, the pressure of the lobes increases and the lobes expand. How fast they expand depends on the density and pressure of the external medium. The highest-pressure phase of the external medium, and thus the most important phase from the point of view of the dynamics, is the X-ray emitting diffuse hot gas. For a long time, it was assumed that powerful sources would expand supersonically, pushing a shock through the external medium. However, X-ray observations show that the internal lobe pressures of powerful FRII sources are often close to the external thermal pressures and not much higher than the external pressures, as would be required for supersonic expansion. The only unambiguously supersonically expanding system known consists of the inner lobes of the low-power radio galaxy Centaurus A which are probably a result of a comparatively recent outburst of the active nucleus.
Host galaxies and environments
These radio sources are almost always found hosted by elliptical galaxies. , Some nearby Seyfert galaxies show weak, small radio jets, but they are neither luminous enough, in radio bands, to be classified as radio-loud nor they are 10-100 times larger than the host spiral galaxies to be comparable to standard radio galaxies. Some observations about the host galaxies of radio-loud quasars and blazars suggest that they are also hosted by elliptical galaxies. There are several possible reasons for this very strong preference for ellipticals. One is that ellipticals generally contain the most massive black holes, and so are capable of powering the most luminous active galaxies (see Eddington luminosity). Another is that ellipticals generally inhabit richer environments, providing a large-scale intergalactic medium to confine the radio source. It may also be that the larger amounts of cold gas in spiral galaxies in some way disrupts or stifles a forming jet. To date there is no compelling single explanation for the observations.
However, since the discovery of the exotic radio galaxy Speca in the year 2011 the number of large-scale radio galaxies hosted in disk/spiral galaxies have been steadily growing. The first such odd-ball radio galaxy was discovered over a decade ago, in 1998 namely J0313-192 in the Abell 428 cluster of galaxies. Possibly due to availability of new sensitive radio and optical all sky surveys the number of such Speca-like objects or Spiral-DRAGNs have increased to up to three dozens. Speca-like such radio galaxies can launch not only giant mega-parsec scale lobes but also be episodic in their nature hosting two pairs of relativistic magnetised plasma lobes. The host spiral galaxies can also be extremely massive spirals and mostly are red in optical colour compared to typical spiral galaxies. It has been speculated that such galaxies may be more common in early Universe when there were less number of ellipticals than spirals.
Unified models
The different types of radio-loud active galaxies are linked by unified models. The key observation that led to the adoption of unified models for powerful radio galaxies and radio-loud quasars was that all quasars appear to be beamed towards us, showing superluminal motion in the cores and bright jets on the side of the source nearest to us (the Laing-Garrington effect:). If this is the case, there must be a population of objects not beamed towards us, and, since we know the lobes are not affected by beaming, they would appear as radio galaxies, provided that the quasar nucleus is obscured when the source is seen side-on. It is now accepted that at least some powerful radio galaxies have 'hidden' quasars, though it is not clear whether all such radio galaxies would be quasars if viewed from the right angle. In a similar way, low-power radio galaxies are a plausible parent population for BL Lac objects.
Uses of radio galaxies
Distant sources
Radio galaxies and radio-loud quasars have been widely used, particularly in the 80s and 90s, to find distant galaxies: by selecting based on radio spectrum and then observing the host galaxy it was possible to find objects at high redshift at modest cost in telescope time. The problem with this method is that hosts of active galaxies may not be typical of galaxies at their redshift. Similarly, radio galaxies have in the past been used to find distant X-ray emitting clusters, but unbiased selection methods are now preferred. The most distant radio galaxy currently known is TGSS J1530+1049, at a redshift of 5.72.
Standard rulers
Some work has been done attempting to use radio galaxies as standard rulers to determine cosmological parameters. This method is fraught with difficulty because a radio galaxy's size depends on both its age and its environment. When a model of the radio source is used, though, methods based on radio galaxies can give good agreement with other cosmological observations.
Effects on environment
Whether or not a radio source is expanding supersonically, it must do work against the external medium in expanding, and so it puts energy into heating and lifting the external plasma. The minimum energy stored in the lobes of a powerful radio source might be 1053 J. The lower limit on the work done on the external medium by such a source is several times this. A good deal of the current interest in radio sources focuses on the effect they must have at the centres of clusters at the present day. Equally interesting is their likely effect on structure formation over cosmological time: it is thought that they may provide a feedback mechanism to slow the formation of the most massive objects.
Terminology
Widely used terminology is awkward now that it is generally accepted that quasars and radio galaxies are the same objects (see above). The acronym DRAGN (for 'Double Radiosource Associated with Galactic Nucleus') was coined by Patrick Leahy in 1993 and is in use. Extragalactic radio source is common but can lead to confusion, since many other extragalactic objects are detected in radio surveys, notably starburst galaxies. Radio-loud active galaxy is unambiguous, and so is often used in this article.
| Physical sciences | Active galactic nucleus | null |
211280 | https://en.wikipedia.org/wiki/Mahi-mahi | Mahi-mahi | The mahi-mahi ( ) or common dolphinfish (Coryphaena hippurus) is a surface-dwelling ray-finned fish found in off-shore temperate, tropical, and subtropical waters worldwide. It is also widely called dorado (not to be confused with Salminus brasiliensis, a freshwater fish) and dolphin (not to be confused with the aquatic mammal dolphin). It is one of two members of the family Coryphaenidae, the other being the pompano dolphinfish. These fish are most commonly found in the waters around the Gulf of Mexico, Costa Rica, Hawaii and the Indian Ocean.
Nomenclature
The name mahi-mahi comes from the Hawaiian language and means 'very strong', through the process of reduplication. By chance in Persian, () means 'fish', but the word is Hawaiian. Though the species is also referred to as the common dolphinfish, the use of dolphin can be misleading as they are not closely related to dolphins; The origin of the name "dolphinfish" is recent and was given to avoid confusion with dolphins, as the traditional name of the fish was also "dolphin". See Coryphaena for the possible etymologies of dolphinfish. In parts of the Pacific and along the English-speaking coast of South Africa, the mahi-mahi is commonly referred to by its name in Spanish, . On the Mediterranean island of Malta, the mahi-mahi is referred to as the . In Indonesian, they are called .
Linnaeus named the genus, derived from the Greek word, , , meaning 'top' or 'apex', in 1758. Synonyms for the species include Coryphaena argyrurus, Coryphaena chrysurus, and Coryphaena dolfyn.
Description
Mahi-mahi have compressed bodies and one very long dorsal fin extending from the head almost to the tail fin. Mature males have distinctive "foreheads"; it grows as the fish matures and often protrudes well above the body proper, which is streamlined by the musculature of the back. This "hump" is a sexually dimorphic feature; females have a rounded head. Their caudal fins and anal fins are sharply concave. They are distinguished by dazzling colors – golden on the sides, and bright blues and greens on the sides and back. The pectoral fins of the mahi-mahi are iridescent blue. The flank is broad and golden.
Out of the water, the fish often change color (giving rise to their Spanish name, , 'golden'), going through several hues before finally fading to a muted yellow-grey upon death.
Mahi-mahi can live for up to five years, although they seldom exceed four. Females are usually smaller than males. Catches typically are and a meter in length. They rarely exceed , and mahi-mahi over are exceptional. Mahi-mahi are among the fastest-growing of fish. They spawn in warm ocean currents throughout much of the year, and their young are commonly found in rafts of Sargassum weeds. Young mahi-mahi migrate past Malta where they are called and Sicily where they are known as or ; there they are fished using nets and floating mats of palm leaves under which they collect.
Mahi-mahi are carnivorous, feeding on flying fish, crabs, squid, mackerel, and other forage fish. They have also been known to eat zooplankton. To pursue such varied pelagic prey, mahi-mahi are fast swimmers, swimming as fast as 50 knots (92.6 km/h, 57.5 mph).
Males and females are sexually mature in their first year, usually by the age of 4–5 months. Spawning can occur at body lengths of . Females may spawn two to three times per year, and produce between 80,000 and 1,000,000 eggs per event. In waters at 28 °C/83 °F, mahi-mahi larvae are found year-round, with greater numbers detected in spring and fall.
Mahi-mahi fish are mostly found in the surface water. Their flesh is grey-white when raw, cooking to an attractive white with a clean, non-fishy flavour.
Relation to humans
Fishing
Recreational fishing
Mahi-mahi are highly sought for sport fishing and commercial purposes. Sport fishermen seek them due to their beauty, size, food quality, and healthy population. Mahi-mahi can be found in the Caribbean Sea, on the west coast of North and South America, the Pacific coast of Costa Rica, the Gulf of Mexico, the Atlantic coast of Florida and West Africa, Indian Ocean, Bay of Bengal, South China Sea and Southeast Asia, Hawaii, Tahiti, and many other places worldwide.
Fishing charters most often look for floating debris and frigatebirds near the edge of the reef in about of water. Mahi-mahi (and many other fish) often swim near debris such as floating wood, five-gallon bucket lids, palm trees and fronds, or sargasso weed lines and around fish buoys. Frigatebirds search for food accompanying the debris or sargasso. Experienced fishing guides can tell what species are likely around the debris by the birds' behavior.
gear is more than adequate when trolling for mahi-mahi. Fly-casters may especially seek frigatebirds to find big mahi-mahis, and then use a bait-and-switch technique. Ballyhoo or a net full of live sardines tossed into the water can excite the mahi-mahis into a feeding frenzy. Hookless teaser lures can have the same effect. After tossing the teasers or live chum, fishermen throw the fly to the feeding mahi-mahi. Once on a line, mahi-mahi are fast, flashy, and acrobatic, with beautiful blue, yellow, green, and even red dots of color.
Commercial fishing
The United States and the Caribbean countries are the primary consumers of this fish, but many European countries are increasing their consumption every year. It is a popular food fish in Australia, usually caught and sold as a byproduct by tuna and swordfish commercial fishing operators. Japan and Hawaii are significant consumers. The Arabian Sea, particularly the coast of Oman, also has mahi-mahi. At first, mahi-mahi were mostly bycatch in the tuna and swordfish longline fishery. Now, they are sought by commercial fishermen on their own merits.
In French Polynesia, fishermen use harpoons, using a specifically designed boat, the , to pursue it, because mahi-mahi do not dive. The is a powerful motorized V-shaped boat, optimized for high agility and speed, and driven with a stick so the pilot can hold his harpoon with his right hand. The method is also practiced by fishermen in the Philippines, especially in the northern province of Batanes, where the harpooning is called .
Environmental and food safety concerns
Depending on how it is caught, mahi-mahi is classed differently by various sustainability rating systems:
The Monterey Bay Aquarium classifies mahi-mahi, when caught in the US Atlantic, as a best choice, the top of its three environmental-impact categories. The aquarium advises to avoid imported mahi-mahi harvested by long line, but rates troll and pole-and-line caught as a good alternative.
The Natural Resources Defense Council classifies mahi-mahi as a "moderate mercury" fish (its second-lowest of four categories), and suggests eating six servings or fewer per month.
The mahi-mahi is also a common vector for ciguatera poisoning. Although a very popular food dish in many parts of the world, there have been reports of ciguatera poisoning from human consumption of this fish. Ciguatera poisoning is caused by the accumulation of toxins (ciguatoxins and maitotoxin) in the flesh of the fish over time. These are produced by Gambierdiscus toxicus which grows together with marine algae, which causes fish like the mahi-mahi to consume them by accident.
Mahi-mahi naturally have high levels of histidine, which is converted to histamine when bacterial growth occurs during improper storage or processing. Subsequent cooking, smoking, or freezing does not eliminate the histamine. This leads to a foodborne illness known as scombroid food poisoning, which also affect other fish such as tuna, mackerel, sardine, anchovy, herring, bluefish, amberjack and marlin. Symptoms are those of histamine intolerance and may include flushed skin, headache, itchiness, blurred vision, abdominal cramps, and diarrhea, and the onset of symptoms is typically 10 to 60 minutes after eating and can last for up to two days. Rarely, breathing problems (like that of allergic asthma) or an irregular heartbeat may occur. Diagnosis is typically based on the symptoms and may be supported by a normal blood tryptase.
| Biology and health sciences | Acanthomorpha | null |
211285 | https://en.wikipedia.org/wiki/Landfill | Landfill | A landfill is a site for the disposal of waste materials. It is the oldest and most common form of waste disposal, although the systematic burial of waste with daily, intermediate and final covers only began in the 1940s. In the past, waste was simply left in piles or thrown into pits (known in archeology as middens).
Landfills take up a lot of land and pose environmental risks. Some landfill sites are used for waste management purposes, such as temporary storage, consolidation and transfer, or for various stages of processing waste material, such as sorting, treatment, or recycling. Unless they are stabilized, landfills may undergo severe shaking or soil liquefaction of the ground during an earthquake. Once full, the area over a landfill site may be reclaimed for other uses.
Operations
Operators of well-run landfills for non-hazardous waste meet predefined specifications by applying techniques to:
confine waste to as small an area as possible
compact waste to reduce volume
They can also cover waste (usually daily) with layers of soil or other types of material such as woodchips and fine particles.
During landfill operations, a scale or weighbridge may weigh waste collection vehicles on arrival and personnel may inspect loads for wastes that do not accord with the landfill's waste-acceptance criteria. Afterward, the waste collection vehicles use the existing road network on their way to the tipping face or working front, where they unload their contents. After loads are deposited, compactors or bulldozers can spread and compact the waste on the working face. Before leaving the landfill boundaries, the waste collection vehicles may pass through a wheel-cleaning facility. If necessary, they return to the weighbridge for re-weighing without their load. The weighing process can assemble statistics on the daily incoming waste tonnage, which databases can retain for record keeping. In addition to trucks, some landfills may have equipment to handle railroad containers. The use of "rail-haul" permits landfills to be located at more remote sites, without the problems associated with many truck trips.
Typically, in the working face, the compacted waste is covered with soil or alternative materials daily. Alternative waste-cover materials include chipped wood or other "green waste", several sprayed-on foam products, chemically "fixed" bio-solids, and temporary blankets. Blankets can be lifted into place at night and then removed the following day prior to waste placement. The space that is occupied daily by the compacted waste and the cover material is called a daily cell. Waste compaction is critical to extending the life of the landfill. Factors such as waste compressibility, waste-layer thickness and the number of passes of the compactor over the waste affect the waste densities.
Sanitary landfill life cycle
The term landfill is usually shorthand for a municipal landfill or sanitary landfill. These facilities were first introduced early in the 20th century, but gained wide use in the 1960s and 1970s, in an effort to eliminate open dumps and other "unsanitary" waste disposal practices. The sanitary landfill is an engineered facility that separates and confines waste. Sanitary landfills are intended as biological reactors (bioreactors) in which microbes will break down complex organic waste into simpler, less toxic compounds over time. These reactors must be designed and operated according to regulatory standards and guidelines (See environmental engineering).
Usually, aerobic decomposition is the first stage by which wastes are broken down in a landfill. These are followed by four stages of anaerobic degradation. Usually, solid organic material in solid phase decays rapidly as larger organic molecules degrade into smaller molecules. These smaller organic molecules begin to dissolve and move to the liquid phase, followed by hydrolysis of these organic molecules, and the hydrolyzed compounds then undergo transformation and volatilization as carbon dioxide (CO2) and methane (CH4), with rest of the waste remaining in solid and liquid phases.
During the early phases, little material volume reaches the leachate, as the biodegradable organic matter of the waste undergoes a rapid decrease in volume. Meanwhile, the leachate's chemical oxygen demand increases with increasing concentrations of the more recalcitrant compounds compared to the more reactive compounds in the leachate. Successful conversion and stabilization of the waste depend on how well microbial populations function in syntrophy, i.e. an interaction of different populations to provide each other's nutritional needs.:
The life cycle of a municipal landfill undergoes five distinct phases:
Initial adjustment (Phase I)
As the waste is placed in the landfill, the void spaces contain high volumes of molecular oxygen (O2). With added and compacted wastes, the O2 content of the landfill bioreactor strata gradually decreases. Microbial populations grow, density increases. Aerobic biodegradation dominates, i.e. the primary electron acceptor is O2.
Transition (Phase II)
The O2 is rapidly degraded by the existing microbial populations. The decreasing O2 leads to less aerobic and more anaerobic conditions in the layers. The primary electron acceptors during transition are nitrates and sulphates since O2 is rapidly displaced by CO2 in the effluent gas.
Acid formation (Phase III)
Hydrolysis of the biodegradable fraction of the solid waste begins in the acid formation phase, which leads to rapid accumulation of volatile fatty acids (VFAs) in the leachate. The increased organic acid content decreases the leachate pH from approximately 7.5 to 5.6. During this phase, the decomposition intermediate compounds like the VFAs contribute much chemical oxygen demand (COD). Long-chain volatile organic acids (VOAs) are converted to acetic acid (C2H4O2), CO2, and hydrogen gas (H2). High concentrations of VFAs increase both the biochemical oxygen demand (BOD) and VOA concentrations, which initiates H2 production by fermentative bacteria, which stimulates the growth of H2-oxidizing bacteria. The H2 generation phase is relatively short because it is complete by the end of the acid formation phase. The increase in the biomass of acidogenic bacteria increases the amount of degradation of the waste material and consuming nutrients. Metals, which are generally more water-soluble at lower pH, may become more mobile during this phase, leading to increasing metal concentrations in the leachate.
Methane fermentation (Phase IV)
The acid formation phase intermediary products (e.g., acetic, propionic, and butyric acids) are converted to CH4 and CO2 by methanogenic microorganisms. As VFAs are metabolized by the methanogens, the landfill water pH returns to neutrality. The leachate's organic strength, expressed as oxygen demand, decreases at a rapid rate with increases in CH4 and CO2 gas production. This is the longest decomposition phase.
Final maturation and stabilization (Phase V)
The rate of microbiological activity slows during the last phase of waste decomposition as the supply of nutrients limits the chemical reactions, e.g. as bioavailable phosphorus becomes increasingly scarce. CH4 production almost completely disappears, with O2 and oxidized species gradually reappearing in the gas wells as O2 permeates downwardly from the troposphere. This transforms the oxidation–reduction potential (ORP) in the leachate toward oxidative processes. The residual organic materials may incrementally be converted to the gas phase, and as organic matter is composted; i.e. the organic matter is converted to humic-like compounds.
Social and environmental impact
Landfills have the potential to cause a number of issues. Infrastructure disruption, such as damage to access roads by heavy vehicles, may occur. Pollution of local roads and watercourses from wheels on vehicles when they leave the landfill can be significant and can be mitigated by wheel washing systems. Pollution of the local environment, such as contamination of groundwater or aquifers or soil contamination may occur, as well.
Leachate
When precipitation falls on open landfills, water percolates through the garbage and becomes contaminated with suspended and dissolved material, forming leachate. If this is not contained it can contaminate groundwater. All modern landfill sites use a combination of impermeable liners several metres thick, geologically stable sites and collection systems to contain and capture this leachate. It can then be treated and evaporated. Once a landfill site is full, it is sealed off to prevent precipitation ingress and new leachate formation. However, liners must have a lifespan, be it several hundred years or more. Eventually, any landfill liner could leak, so the ground around landfills must be tested for leachate to prevent pollutants from contaminating groundwater.
Decomposition gases
Rotting food and other decaying organic waste create decomposition gases, especially CO2 and CH4 from aerobic and anaerobic decomposition, respectively. Both processes occur simultaneously in different parts of a landfill. In addition to available O2, the fraction of gas constituents will vary, depending on the age of landfill, type of waste, moisture content and other factors. For example, the maximum amount of landfill gas produced can be illustrated a simplified net reaction of diethyl oxalate that accounts for these simultaneous reactions:
4 C6H10O4 + 6 H2O → 13 CH4 + 11 CO2
On average, about half of the volumetric concentration of landfill gas is CH4 and slightly less than half is CO2. The gas also contains about 5% molecular nitrogen (N2), less than 1% hydrogen sulfide (H2S), and a low concentration of non-methane organic compounds (NMOC), about 2700 ppmv.
Landfill gases can seep out of the landfill and into the surrounding air and soil. Methane is a greenhouse gas, and is flammable and potentially explosive at certain concentrations, which makes it perfect for burning to generate electricity cleanly. Since decomposing plant matter and food waste only release carbon that has been captured from the atmosphere through photosynthesis, no new carbon enters the carbon cycle and the atmospheric concentration of CO2 is not affected. Carbon dioxide traps heat in the atmosphere, contributing to climate change. In properly managed landfills, gas is collected and flared or recovered for landfill gas utilization.
Vectors
Poorly run landfills may become nuisances because of vectors such as rats and flies which can spread infectious diseases. The occurrence of such vectors can be mitigated through the use of daily cover.
Other nuisances
Other potential issues include wildlife disruption due to occupation of habitat and animal health disruption caused by consuming waste from landfills, dust, odor, noise pollution, and reduced local property values.
Landfill gas
Gases are produced in landfills due to the anaerobic digestion by microbes. In a properly managed landfill, this gas is collected and used. Its uses range from simple flaring to the landfill gas utilization and generation of electricity. Landfill gas monitoring alerts workers to the presence of a build-up of gases to a harmful level. In some countries, landfill gas recovery is extensive; in the United States, for example, more than 850 landfills have active landfill gas recovery systems.
Solar landfill
A Solar landfill is a repurposed used landfill that is converted to a solar array solar farm.
Regional practice
Canada
Landfills in Canada are regulated by provincial environmental agencies and environmental protection legislation.
Older facilities tend to fall under current standards and are monitored for leaching. Some former locations have been converted to parkland.
European Union
In the European Union, individual states are obliged to enact legislation to comply with the requirements and obligations of the European Landfill Directive.
The majority of EU member states have laws banning or severely restricting the disposal of household trash via landfills.
India
Landfilling is currently the major method of municipal waste disposal in India. India also has Asia's largest dumping ground in Deonar, Mumbai. However, issues frequently arise due to the alarming growth rate of landfills and poor management by authorities. On and under surface fires have been commonly seen in the Indian landfills over the last few years.
United Kingdom
Landfilling practices in the UK have had to change in recent years to meet the challenges of the European Landfill Directive. The UK now imposes landfill tax upon biodegradable waste which is put into landfills. In addition to this the Landfill Allowance Trading Scheme has been established for local authorities to trade landfill quotas in England. A different system operates in Wales where authorities cannot 'trade' amongst themselves, but have allowances known as the Landfill Allowance Scheme.
United States
U.S. landfills are regulated by each state's environmental agency, which establishes minimum guidelines; however, none of these standards may fall below those set by the United States Environmental Protection Agency (EPA).
Permitting a landfill generally takes between five and seven years, costs millions of dollars and requires rigorous siting, engineering and environmental studies and demonstrations to ensure local environmental and safety concerns are satisfied.
Types
Municipal solid waste: takes in household waste and nonhazardous material. Included in this type of landfill is a Bioreactor Landfill that specifically degrades organic material.
Industrial waste: for commercial and industrial waste. Other related landfills include Construction and Demolition Debris Landfills and Coal Combustion Residual Landfills.
Hazardous waste or PCB waste: Polychlorinated Biphenyl (PCB) landfills that are monitored in the United States by the Toxic Substances Control Act of 1976 (TSCA).
Microbial topics
The status of a landfill's microbial community may determine its digestive efficiency.
Bacteria that digest plastic have been found in landfills.
Reclaiming materials
One can treat landfills as a viable and abundant source of materials and energy. In the developing world, waste pickers often scavenge for still-usable materials. In commercial contexts, companies have also discovered landfill sites, and many have begun harvesting materials and energy. Well-known examples include gas-recovery facilities.
Other commercial facilities include waste incinerators which have built-in material recovery. This material recovery is possible through the use of filters (electro filter, active-carbon and potassium filter, quench, HCl-washer, SO2-washer, bottom ash-grating, etc.).
Alternatives
In addition to waste reduction and recycling strategies, there are various alternatives to landfills, including waste-to-energy incineration, anaerobic digestion, composting, mechanical biological treatment, pyrolysis and plasma arc gasification. Depending on local economics and incentives, these can be made more financially attractive than landfills.
The goal of the zero waste concept is to minimize landfill volume.
Restrictions
Countries including Germany, Austria, Sweden, Denmark, Belgium, the Netherlands, and Switzerland, have banned the disposal of untreated waste in landfills. In these countries, only certain hazardous wastes, fly ashes from incineration or the stabilized output of mechanical biological treatment plants may still be deposited.
| Technology | Basics_6 | null |
211324 | https://en.wikipedia.org/wiki/Turaco | Turaco | The turacos make up the bird family Musophagidae ( "banana-eaters"), which includes plantain-eaters and go-away-birds. In southern Africa both turacos and go-away-birds are commonly known as loeries. They are semi-zygodactylous: the fourth (outer) toe can be switched back and forth. The second and third toes, which always point forward, are conjoined in some species. Musophagids often have prominent crests and long tails; the turacos are noted for peculiar and unique pigments giving them their bright green and red feathers.
Traditionally, this group has been allied with the cuckoos in the order Cuculiformes, but the Sibley-Ahlquist taxonomy raises this group to a full order Musophagiformes. They have been proposed to link the hoatzin to the other living birds, but this was later disputed. Recent genetic analyses have strongly supported the order ranking of Musophagiformes.
Musophagidae is one of very few bird families endemic to Africa, one other being the mousebirds, Colliidae. All species are frugivorous, but they also eat leaves, buds, and flowers. Figs are an important part of their diet. They have rounded wings and long tails and strong legs, making them poor fliers, but good runners.
Turacos are medium-sized arboreal birds endemic to sub-Saharan Africa, where they live in forests, woodland and savanna. Turacos can occasionally be found outside of their native range as escapes from captivity.
They are gregarious, non-migratory birds which move in family groups of up to 10. Many species are noisy, with the go-away-birds being especially noted for their piercing alarm calls, which alert other fauna to the presence of predators; their common name is onomatopoeia of this. Musophagids build large stick nests in trees, and lay 2 or 3 eggs. The young are born with thick down and open, or nearly-open, eyes.
Morphology
Most turacos are medium-sized birds – an exception being the large great blue turaco – with long tails and short, rounded wings. They range in length from . Their flight is weak, but they are strong climbers and are able to move nimbly on branches and through vegetation. Juveniles have claws on the wings that help them climb. They have a unique foot arrangement, where the fourth toe can be brought around to the back of the foot where it almost touches the first toe, or brought around so that it is near the second and third. In spite of this flexibility the toe is actually usually held at right angles to the axis of the foot.
The plumage of go-away-birds and plantain-eaters is mainly grey and white. The turacos on the other hand are brightly coloured birds, usually blue, green or purple. The green colour in turacos comes from turacoverdin, the only true green pigment in birds known to date. Other "greens" in bird colors result from a yellow pigment such as some carotenoid, combined with the prismatic physical structure of the feather itself which scatters the light in a particular way and giving a blue colour.
Turaco wings contain the red pigment turacin, unlike in other birds where red colour is due to carotenoids. Both pigments are derived from porphyrins and only known from the Musophagidae into the 21st century, but especially the little-researched turacoverdin might have relatives in other birds. The incidence of turacoverdin in relation to habitat is of interest to scientists, being present in forest species but absent in savanna- and acacia-living species.
Little is known about the longevity of wild turacos, but in captivity they are proving to be exceptionally long-lived, easily living to 30 years in captivity. A bird in the Cotswold Wildlife Park collection in England approached its 37th year.
Evolution and systematics
The fossil genus Veflintornis is known from the Middle Miocene of Grive-Saint-Alban (France). It was established as Apopempsis by Pierce Brodkorb in 1971, but this is pre-empted by Schenkling's 1903 use of the name for some beetles. "Apopempsis" africanus (Early Miocene of Kenya) might also belong there.
Further fossil material of putative musophagids was found in Egypt as well as in Late Oligocene deposits at Gaimersheim in Germany and Middle Miocene deposits at Grive-Saint-Alban and Vieux-Collonges (each in France). While it is not entirely certain that these fossils are indeed of turacos, it nonetheless appears as if the family evolved in the Oligocene of central Europe or perhaps northern Africa, and later on shifted its distribution southwards. The climate of those European regions during the late Paleogene was not too dissimilar to that of (sub)tropical Africa today; the Saharan desert was not yet present and the distance across the Mediterranean was not much more than what it is today. Thus such a move south may well have been a very slow and gradual shifting of a large and continuous range.
The Early Eocene Promusophaga was initially believed to be the oldest record of the turacos; it was eventually reconsidered a distant relative of the ostrich and is now in the ratite family Lithornithidae. Filholornis from the Late Eocene or Early Oligocene of France is occasionally considered a musophagid, but its relationships have always been disputed. It is not often considered a turaco in more recent times and has been synonymised with the presumed gruiform Talantatos, though it is not certain whether this will become widely accepted.
The phylogenetic analysis conducted by Field & Hsiang (2018) indicated that Eocene (Wasatchian) species Foro panarium known from the Fossil Butte Member of the Green River Formation (Wyoming, United States) was a stem-turaco.
Phylogeny
The IOC World Bird List (version 10.1) recognises 23 species of turaco in six genera. However, a phylogenetic analysis by Perktaş et al (2020) found genus Tauraco polyphyletic and a revised classification has been proposed based on molecular, morphological and biogeographic analysis. This study recognised 33 species-level taxa in seven genera corresponding to the major clades. The following phylogenetic tree is based on this proposal and uses their proposed genus and species names.
| Biology and health sciences | Cuculiformes | null |
211386 | https://en.wikipedia.org/wiki/Scapula | Scapula | The scapula (: scapulae or scapulas), also known as the shoulder blade, is the bone that connects the humerus (upper arm bone) with the clavicle (collar bone). Like their connected bones, the scapulae are paired, with each scapula on either side of the body being roughly a mirror image of the other. The name derives from the Classical Latin word for trowel or small shovel, which it was thought to resemble.
In compound terms, the prefix omo- is used for the shoulder blade in medical terminology. This prefix is derived from ὦμος (ōmos), the Ancient Greek word for shoulder, and is cognate with the Latin , which in Latin signifies either the shoulder or the upper arm bone.
The scapula forms the back of the shoulder girdle. In humans, it is a flat bone, roughly triangular in shape, placed on a posterolateral aspect of the thoracic cage.
Structure
The scapula is a thick, flat bone lying on the thoracic wall that provides an attachment for three groups of muscles: intrinsic, extrinsic, and stabilizing and rotating muscles.
The intrinsic muscles of the scapula include the muscles of the rotator cuff(SITS muscle)—the subscapularis, supraspinatus, infraspinatus and teres minor. These muscles attach to the surface of the scapula and are responsible for the internal and external rotation of the shoulder joint, along with humeral abduction.
The extrinsic muscles include the biceps, triceps, and deltoid muscles and attach to the coracoid process and supraglenoid tubercle of the scapula, infraglenoid tubercle of the scapula, and spine of the scapula. These muscles are responsible for several actions of the glenohumeral joint.
The third group, which is mainly responsible for stabilization and rotation of the scapula, consists of the trapezius, serratus anterior, levator scapulae, and rhomboid muscles. These attach to the medial, superior, and inferior borders of the scapula.
The head, processes, and the thickened parts of the bone contain cancellous tissue; the rest consists of a thin layer of compact tissue.
The central part of the supraspinatus fossa and the upper part of the infraspinatous fossa, but especially the former, are usually so thin in humans as to be semitransparent; occasionally the bone is found wanting in this situation, and the adjacent muscles are separated only by fibrous tissue. The scapula has two surfaces, three borders, three angles, and three processes.
Surfaces
Front or subscapular fossa
The front of the scapula (also known as the costal or ventral surface) has a broad concavity called the subscapular fossa, to which the subscapularis muscle attaches. The medial two-thirds of the fossa have 3 longitudinal oblique ridges, and another thick ridge adjoins the lateral border; they run outward and upward. The ridges give attachment to the tendinous insertions, and the surfaces between them to the fleshy fibers, of the subscapularis muscle. The lateral third of the fossa is smooth and covered by the fibers of this muscle.
At the upper part of the fossa is a transverse depression, where the bone appears to be bent on itself along a line at right angles to and passing through the center of the glenoid cavity, forming a considerable angle, called the subscapular angle; this gives greater strength to the body of the bone by its arched form, while the summit of the arch serves to support the spine and acromion.
The costal surface superior of the scapula is the origin of 1st digitation for the serratus anterior origin.
Back
The back of the scapula (also called the dorsal or posterior surface) is arched from above downward, and is subdivided into two unequal parts by the spine of the scapula. The portion above the spine is called the supraspinous fossa, and that below it the infraspinous fossa. The two fossae are connected by the spinoglenoid notch, situated lateral to the root of the spine.
The supraspinous fossa, above the spine of scapula, is concave, smooth, and broader at its vertebral than at its humeral end; its medial two-thirds give origin to the Supraspinatus. At its lateral surface resides the spinoglenoid fossa which is situated by the medial margin of the glenoid. The spinoglenoid fossa houses the suprascapular canal which forms a connecting passage between the suprascapular notch and the spinoglenoid notch conveying the suprascapular nerve and vessels.
The infraspinous fossa is much larger than the preceding; toward its vertebral margin a shallow concavity is seen at its upper part; its center presents a prominent convexity, while near the axillary border is a deep groove which runs from the upper toward the lower part. The medial two-thirds of the fossa give origin to the Infraspinatus; the lateral third is covered by this muscle.
There is a ridge on the outer part of the back of the scapula. This runs from the lower part of the glenoid cavity, downward and backward to the vertebral border, about 2.5 cm above the inferior angle. Attached to the ridge is a fibrous septum, which separates the infraspinatus muscle from the Teres major and Teres minor muscles. The upper two-thirds of the surface between the ridge and the axillary border is narrow, and is crossed near its center by a groove for the scapular circumflex vessels; the Teres minor attaches here.
The broad and narrow portions above alluded to are separated by an oblique line, which runs from the axillary border, downward and backward, to meet the elevated ridge: to it is attached a fibrous septum which separates the Teres muscles from each other.
Its lower third presents a broader, somewhat triangular surface, the inferior angle of the scapula, which gives origin to the Teres major, and over which the Latissimus dorsi glides; frequently the latter muscle takes origin by a few fibers from this part.
Side
The acromion forms the summit of the shoulder, and is a large, somewhat triangular or oblong process, flattened from behind forward, projecting at first laterally, and then curving forward and upward, so as to overhang the glenoid cavity.
Angles
There are 3 angles:
The superior angle of the scapula or medial angle, is covered by the trapezius muscle. This angle is formed by the junction of the superior and medial borders of the scapula. The superior angle is located at the approximate level of the second thoracic vertebra. The superior angle of the scapula is thin, smooth, rounded, and inclined somewhat lateralward, and gives attachment to a few fibers of the levator scapulae muscle.
The inferior angle of the scapula is the lowest part of the scapula and is covered by the latissimus dorsi muscle. It moves forwards round the chest when the arm is abducted. The inferior angle is formed by the union of the medial and lateral borders of the scapula. It is thick and rough and its posterior or back surface affords attachment to the teres major and often to a few fibers of the latissimus dorsi. The anatomical plane that passes vertically through the inferior angle is named the scapular line.
The lateral angle of the scapula or glenoid angle, also known as the head of the scapula, is the thickest part of the scapula. It is broad and bears the glenoid fossa on its articular surface which is directed forward, laterally and slightly upwards, and articulates with the head of the humerus. The inferior angle is broader below than above and its vertical diameter is the longest. The surface is covered with cartilage in the fresh state; and its margins, slightly raised, give attachment to a fibrocartilaginous structure, the glenoidal labrum, which deepens the cavity. At its apex is a slight elevation, the supraglenoid tuberosity, to which the long head of the biceps brachii is attached.
The anatomic neck of the scapula is the slightly constricted portion which surrounds the head and is more distinct below and behind than above and in front. The surgical neck of the scapula passes directly medial to the base of the coracoid process.
Borders
There are three borders of the scapula:
The superior border is the shortest and thinnest; it is concave, and extends from the superior angle to the base of the coracoid process. It is referred to as the cranial border in animals.
At its lateral part is a deep, semicircular notch, the scapular notch, formed partly by the base of the coracoid process. This notch is converted into a foramen by the superior transverse scapular ligament, and serves for the passage of the suprascapular nerve; sometimes the ligament is ossified.
The adjacent part of the superior border affords attachment to the omohyoideus.
The axillary border (or "lateral border") is the thickest of the three. It begins above at the lower margin of the glenoid cavity, and inclines obliquely downward and backward to the inferior angle. It is referred to as the caudal border in animals.
It begins above at the lower margin of the glenoid cavity, and inclines obliquely downward and backward to the inferior angle.
Immediately below the glenoid cavity is a rough impression, the infraglenoid tuberosity, about . in length, which gives origin to the long head of the triceps brachii; in front of this is a longitudinal groove, which extends as far as the lower third of this border, and affords origin to part of the subscapularis.
The inferior third is thin and sharp, and serves for the attachment of a few fibers of the teres major behind, and of the subscapularis in front.
The medial border (also called the vertebral border or medial margin) is the longest of the three borders, and extends from the superior angle to the inferior angle. In animals it is referred to as the dorsal border.
Four muscles attach to the medial border. Serratus anterior has a long attachment on the anterior lip. Three muscles insert along the posterior lip, the levator scapulae (uppermost), rhomboid minor (middle), and to the rhomboid major (lower middle).
Development
The scapula is ossified from 7 or more centers: one for the body, two for the coracoid process, two for the acromion, one for the vertebral border, and one for the inferior angle. Ossification of the body begins about the second month of fetal life, by an irregular quadrilateral plate of bone forming, immediately behind the glenoid cavity. This plate extends to form the chief part of the bone, the scapular spine growing up from its dorsal surface about the third month. Ossification starts as membranous ossification before birth. After birth, the cartilaginous components would undergo endochondral ossification. The larger part of the scapula undergoes membranous ossification. Some of the outer parts of the scapula are cartilaginous at birth, and would therefore undergo endochondral ossification.
At birth, a large part of the scapula is osseous, but the glenoid cavity, the coracoid process, the acromion, the vertebral border and the inferior angle are cartilaginous. From the 15th to the 18th month after birth, ossification takes place in the middle of the coracoid process, which as a rule becomes joined with the rest of the bone about the 15th year.
Between the 14th and 20th years, the remaining parts ossify in quick succession, and usually in the following order: first, in the root of the coracoid process, in the form of a broad scale; secondly, near the base of the acromion; thirdly, in the inferior angle and contiguous part of the vertebral border; fourthly, near the outer end of the acromion; fifthly, in the vertebral border.
The base of the acromion is formed by an extension from the spine; the two nuclei of the acromion unite, and then join with the extension from the spine. The upper third of the glenoid cavity is ossified from a separate center (sub coracoid), which appears between the 10th and 11th years and joins between the 16th and the 18th years. Further, an epiphysial plate appears for the lower part of the glenoid cavity, and the tip of the coracoid process frequently has a separate nucleus. These various epiphyses are joined to the bone by the 25th year.
Failure of bony union between the acromion and spine sometimes occurs (see os acromiale), the junction being effected by fibrous tissue, or by an imperfect articulation; in some cases of supposed fracture of the acromion with ligamentous union, it is probable that the detached segment was never united to the rest of the bone.
Function
The following muscles attach to the scapula:
Movements
Movements of the scapula are brought about by the scapular muscles. The scapula can perform six actions:
Elevation: upper trapezius and levator scapulae
Depression: lower trapezius
Retraction (adduction): rhomboids and middle trapezius
Protraction (abduction): serratus anterior
Upward rotation: upper and lower trapezius, serratus anterior
Downward rotation: rhomboids, Levator Scapulae, and Pec Minor
Clinical significance
Scapular fractures
Because of its sturdy structure and protected location, fractures of the scapula are uncommon. When they do occur, they are an indication that severe chest trauma has occurred.
Scapular fractures involving the neck of the scapula have two patterns. One (rare) type of fracture is through the anatomical neck of the scapula. The other more common type of fracture is through the surgical neck of the scapula. The surgical neck exits medial to the coracoid process.
An abnormally protruding inferior angle of the scapula is known as a winged scapula and can be caused by paralysis of the serratus anterior muscle. In this condition the sides of the scapula nearest the spine are positioned outward and backward. The appearance of the upper back is said to be wing-like. In addition, any condition causing weakness of the serratus anterior muscle may cause scapular "winging".
Scapular dyskenesis
The scapula plays an important role in shoulder impingement syndrome.
Abnormal scapular function is called scapular dyskinesis. The scapula performs elevation of the acromion process during a throwing or serving motion, in order to avoid impingement of the rotator cuff tendons. If the scapula fails to properly elevate the acromion, impingement may occur during the cocking and acceleration phase of an overhead activity. The two muscles most commonly inhibited during this first part of an overhead motion are the serratus anterior and the lower trapezius. These two muscles act as a force couple within the glenohumeral joint to properly elevate the acromion process, and if a muscle imbalance exists, shoulder impingement may develop.
Other conditions associated with scapular dyskenesis include thoracic outlet syndrome and the related pectoralis minor syndrome.
Etymology
The name scapula as synonym of shoulder blade is of Latin origin. It is commonly used in medical English and is part of the current official Latin nomenclature, Terminologia Anatomica. Shoulder blade is the colloquial name for this bone.
In other animals
In fish, the scapular blade is a structure attached to the upper surface of the articulation of the pectoral fin, and is accompanied by a similar coracoid plate on the lower surface. Although sturdy in cartilagenous fish, both plates are generally small in most other fish, and may be partially cartilagenous, or consist of multiple bony elements.
In the early tetrapods, these two structures respectively became the scapula and a bone referred to as the procoracoid (commonly called simply the "coracoid", but not homologous with the mammalian structure of that name). In amphibians and reptiles (birds included), these two bones are distinct, but together form a single structure bearing many of the muscle attachments for the forelimb. In such animals, the scapula is usually a relatively simple plate, lacking the projections and spine that it possesses in mammals. However, the detailed structure of these bones varies considerably in living groups. For example, in frogs, the procoracoid bones may be braced together at the animal's underside to absorb the shock of landing, while in turtles, the combined structure forms a Y-shape in order to allow the scapula to retain a connection to the clavicle (which is part of the shell). In birds, the procoracoids help to brace the wing against the top of the sternum.
In the fossil therapsids, a third bone, the true coracoid, formed just behind the procoracoid. The resulting three-boned structure is still seen in modern monotremes, but in all other living mammals, the procoracoid has disappeared, and the coracoid bone has fused with the scapula, to become the coracoid process. These changes are associated with the upright gait of mammals, compared with the more sprawling limb arrangement of reptiles and amphibians; the muscles formerly attached to the procoracoid are no longer required. The altered musculature is also responsible for the alteration in the shape of the rest of the scapula; the forward margin of the original bone became the spine and acromion, from which the main shelf of the shoulder blade arises as a new structure.
In dinosaurs
In dinosaurs the main bones of the pectoral girdle were the scapula (shoulder blade) and the coracoid, both of which directly articulated with the clavicle. The clavicle was present in saurischian dinosaurs but largely absent in ornithischian dinosaurs. The place on the scapula where it articulated with the humerus (upper bone of the forelimb) is called the glenoid. The scapula serves as the attachment site for a dinosaur's back and forelimb muscles.
Gallery
| Biology and health sciences | Skeletal system | Biology |
211397 | https://en.wikipedia.org/wiki/Toe | Toe | Toes are the digits of the foot of a tetrapod. Animal species such as cats that walk on their toes are described as being digitigrade. Humans, and other animals that walk on the soles of their feet, are described as being plantigrade; unguligrade animals are those that walk on hooves at the tips of their toes.
Structure
There are normally five toes present on each human foot. Each toe consists of three phalanx bones, the proximal, middle, and distal, with the exception of the big toe (). For a minority of people, the little toe also is missing a middle bone. The hallux only contains two phalanx bones, the proximal and distal. The joints between each phalanx are the interphalangeal joints. The proximal phalanx bone of each toe articulates with the metatarsal bone of the foot at the metatarsophalangeal joint. Each toe is surrounded by skin, and present on all five toes is a toenail.
The toes are, from medial to lateral:
the first toe, also known as the hallux ("big toe", "great toe", "thumb toe"), the innermost toe;
the second toe, ("index toe", "pointer toe");
the third toe, ("middle toe");
the fourth toe, ("fore toe", "ring toe");
the fifth toe, ("baby toe", "little toe", "pinky toe", "small toe"), the outermost toe.
Muscles
Toe movement is generally flexion and extension (movement toward the sole or the back of the foot, resp.) via muscular tendons that attach to the toes on the anterior and superior surfaces of the phalanx bones.
With the exception of the hallux, toe movement is generally governed by action of the flexor digitorum brevis and extensor digitorum brevis muscles. These attach to the sides of the bones, making it impossible to move individual toes independently. Muscles between the toes on their top and bottom also help to abduct and adduct the toes. The hallux and little toe have unique muscles:
The hallux is primarily flexed by the flexor hallucis longus muscle, located in the deep posterior of the lower leg, via the flexor hallucis longus tendon. Additional flexion control is provided by the flexor hallucis brevis. It is extended by the abductor hallucis muscle and the adductor hallucis muscle.
The little toe has a separate set of control muscles and tendon attachments, the flexor and abductor digiti minimi. Numerous other foot muscles contribute to fine motor control of the foot. The connective tendons between the minor toes account for the inability to actuate individual toes.
Blood supply
The toes receive blood from the digital branches of the plantar metatarsal arteries and drain blood into the dorsal venous arch of the foot.
Nerve supply
Sensation to the skin of the toes is provided by five nerves. The superficial fibular nerve supplies sensation to the top of the toes, except between the hallux and second toe, which is supplied by the deep fibular nerve, and the outer surface of the fifth toe, supplied by the sural nerve. Sensation to the bottom of the toes is supplied by the medial plantar nerve, which supplies sensation to the great toe and inner three-and-a-half toes, and the lateral plantar nerve, which supplies sensation to the little toe and half of the sensation of the fourth toe.
In humans, the hallux is usually longer than the second toe, but in some individuals, it may not be the longest toe. There is an inherited trait in humans, where the dominant gene causes a longer second toe ("Morton's toe" or "Greek foot") while the homozygous recessive genotype presents with the more common trait: a longer hallux. People with the rare genetic disease fibrodysplasia ossificans progressiva characteristically have a short hallux which appears to turn inward, or medially, in relation to the foot.
Variation
Humans usually have five toes on each foot. When more than five toes are present, this is known as polydactyly. Other variants may include syndactyly or arachnodactyly. Forefoot shape, including toe shape, exhibits significant variation among people; these differences can be measured and have been statistically correlated with ethnicity. Such deviations may affect comfort and fit for various shoe types. Research conducted for the U.S. Army indicated that larger feet may still have smaller arches, toe length, and toe-breadth.
Function
The human foot consists of multiple bones and soft tissues which support the weight of the upright human. Specifically, the toes assist the human while walking, providing balance, weight-bearing, and thrust during gait.
Clinical significance
A sprain or strain to the small interphalangeal joints of the toe is commonly called a stubbed toe. A sprain or strain where the toe joins to the foot is called turf toe.
Long-term use of improperly sized shoes can cause misalignment of toes, as well as other orthopedic problems.
Morton's neuroma commonly results in pain and numbness between the third and fourth toes of the sufferer, due to it affecting the nerve between the third and fourth metatarsal bones.
The big toe is also the most common locus of ingrown nails, and its proximal phalanx joint is the most common locus for gout attacks.
Deformity
Deformities of the foot include hammer toe, trigger toe, and claw toe. Hammer toe can be described as an abnormal contraction or “buckling” of a toe. This is done by a partial or complete dislocation of one of the joints, which form the toes. Since the toes are deformed further, these may press against a shoe and cause pain. Deformities of the foot can also be caused by rheumatoid arthritis and diabetes mellitus. Deformities may predispose to ulcers and pain, especially when shoe-wearing.
A common problem involving the big toe is the formation of bunions. These are structural deformities of the bones and the joint between the foot and big toe, and may be painful. Similar deformity involving the fifth toe is described as tailor's bunion or bunionette.
In polydactyly (which can also affect the fingers) one or more extra toes are present.
In reconstruction
A favourable option for the reconstruction of missing adjacent fingers/multiple digit amputations, i.e. such as a metacarpal hand reconstruction, is to have a combined second and third toe transplantation. Third and fourth toe transplantation to the hand in replacing lost fingers is also a viable option.
History
Etymology
The Old English term for toe is (plural ). This is a contraction of , and derives from Proto-Germanic (cognates: Old Norse , Old Frisian , Middle Dutch , Dutch (perhaps originally a plural), Old High German , German ), perhaps originally meaning 'fingers' as well (many Indo-European languages use one word to mean both 'fingers' and 'toes', e.g. digit), and thus from PIE root — 'to show'.
In classical Latin, , , and , with genitive and , are used to refer to the big toe. The form (genitive, ) currently in use is however a blend word of the aforementioned forms. Compare pollex, the equivalent term for the thumb.
Evolution
Haeckel traces the standard vertebrate five-toed schema from fish fins via amphibian ancestors.
Other animals
In birds with anisodactyl or heterodactyl feet, the hallux is opposed or directed backwards and allows for grasping and perching.
While the thumb is often mentioned as one of the signature characteristics in humans, this manual digit remains partially primitive and is actually present in all primates. In humans, the most derived digital feature is the hallux.
| Biology and health sciences | Human anatomy | Health |
211411 | https://en.wikipedia.org/wiki/Heteroptera | Heteroptera | The Heteroptera are a group of about 40,000 species of insects in the order Hemiptera. They are sometimes called "true bugs", though that name more commonly refers to the Hemiptera as a whole. "Typical bugs" might be used as a more unequivocal alternative, since the heteropterans are most consistently and universally termed "bugs" among the Hemiptera. "Heteroptera" is Greek for "different wings": most species have forewings with both membranous and hardened portions (called hemelytra); members of the primitive sub-group Enicocephalomorpha have completely membranous wings.
The name "Heteroptera" is used in two very different ways in modern classifications. In Linnean nomenclature, it commonly appears as a suborder within the order Hemiptera, where it can be paraphyletic or monophyletic depending on its delimitation. In phylogenetic nomenclature, it is used as an unranked clade within the Prosorrhyncha clade, which in turn is in the Hemiptera clade. This results from the realization that the Coleorrhyncha are just "living fossil" relatives of the traditional Heteroptera, close enough to them to be united with that group.
The infraorders Leptopodomorpha, Gerromorpha, and Nepomorpha, comprise a significant component of the world's aquatic and semiaquatic insects. There are 23 families, 343 genera and 4,810 species group taxa within these three infraorders. Most of the remaining groups that are common and familiar are in the Cimicomorpha and Pentatomomorpha.
Classification
The use of the name "Heteroptera" has had the rank of order, dating back to 1810 by Pierre André Latreille. Only recently has it been relegated to a subsidiary rank within a larger definition of Hemiptera, so many reference works still include it as an order. Whether to continue treating it as a suborder is still a subject of some controversy, as is whether the name itself should ever be used, although three basic approaches ranging from abolishing it entirely to maintaining the taxonomy with a slight change in systematics is proposed, two of which (but not the traditional one) agree with the phylogeny. The competing classifications call for a preference for two suborders versus one when the "living fossil" family Peloridiidae is taken into consideration:
In one revised classification proposed in 1995, the name of the suborder is Prosorrhyncha, and Heteroptera is a rankless subgroup within it. The only difference between Heteroptera and Prosorrhyncha is that the latter includes the family Peloridiidae, which is a tiny relictual group that is in its own monotypic superfamily and infraorder. In other words, the Heteroptera and Prosorrhyncha sensu Sorensen et al. are identical except that Prosorrhyncha contains one additional infraorder, called Peloridiomorpha (comprising only 13 small genera). The ongoing conflict between traditional, Linnaean classifications and nontraditional classifications is exemplified by the problem inherent in continued usage of the name Heteroptera when it no longer can be matched to any standard Linnaean rank (as it falls below suborder but above infraorder). If this classification succeeds, then the "Heteroptera" grouping may be discarded, but in that case it is likely that no ranks will be used at all according to the standards of phylogenetic nomenclature.
In the traditional classification, the Peloridiidae are retained as their own suborder, called Coleorrhyncha; "Heteroptera" is treated the same. Functionally, the only difference between this classification and the preceding is that the former uses the name Prosorrhyncha to refer to a particular clade, while the traditional approach divides this into the paraphyletic Heteroptera and the monophyletic Coleorrhyncha. Many believe it is preferable to use only one name because the characteristics of the two traditional suborders are too closely related to be treated as separate.
Alternatively, the modified approach of placing Coleorrhyncha 'within' the Heteroptera can be used. Indeed, as that solution preserves the well-known Heteroptera at the taxonomic rank they traditionally hold while making them a good monophyletic group, it seems preferable to the paraphyletic "Heteroptera" used in older works. In that case, the "core" Heteroptera could be considered a section – as yet unnamed, mainly because the Prosorrhyncha were proposed earlier – within the "expanded" Heteroptera, or the latter could simply be described as consisting of a basal "living fossil" lineage and a more apomorphic main radiation. Whether the name "Coleorrhyncha" is to be retained for the basal lineage or whether the more consistent "Peloridiomorpha" is used instead is a matter of taste, as described below.
Separate from the question of the actual "closeness" of Heteroptera and Coleorrhyncha is the potential disruption to traditional construction of names; there seems to be reluctance among hemipterists to abandon the use of "Heteroptera". This can be seen by the name itself, as it is a violation of convention to use the ending "-ptera" for any rank above genus other than an order – though since it is a convention rather than a mandatory rule of Linnean nomenclature, taxonomists are technically free to violate it (which is why, for example, not all insect orders end in "-ptera", e.g., Odonata). However, in most cases when such conventions are violated, it does not create an internal conflict as in the present case (that is, the order Hemiptera has a suborder named Heteroptera, which is an internal conflict). At least some hemipterists argue that the name Heteroptera should be dropped entirely to eliminate this internal conflict, though the third possibility offers a workaround. In that case, to achieve full consistency of names "Coleorrhyncha" would probably be dropped in favor of "Peloridiomorpha".
| Biology and health sciences | Hemiptera (true bugs) | Animals |
211426 | https://en.wikipedia.org/wiki/Gnathostomata | Gnathostomata | Gnathostomata (; from Ancient Greek: () 'jaw' + () 'mouth') are the jawed vertebrates. Gnathostome diversity comprises roughly 60,000 species, which accounts for 99% of all living vertebrates, including humans. Most gnathostomes have retained ancestral traits like true teeth, a stomach, and paired appendages (pectoral and pelvic fins, arms, legs, wings, etc.). Other traits are elastin, a horizontal semicircular canal of the inner ear, myelin sheaths of neurons, and an adaptive immune system which has discrete lymphoid organs (spleen and thymus), and uses V(D)J recombination to create antigen recognition sites, rather than using genetic recombination in the variable lymphocyte receptor gene.
It is now assumed that Gnathostomata evolved from ancestors that already possessed a pair of both pectoral and pelvic fins. Until recently these ancestors, known as antiarchs, were thought to have lacked pectoral or pelvic fins. In addition to this, some placoderms (extinct fish with bony plates) were shown to have a third pair of paired appendages, that had been modified to claspers in males and basal plates in females—a pattern not seen in any other vertebrate group.
The Osteostraci (bony armored jawless fish) are generally considered the sister taxon of Gnathostomata.
Jaw development in vertebrates is likely a product of the supporting gill arches. This development would help push water into the mouth by the movement of the jaw, so that it would pass over the gills for gas exchange. The repetitive use of the newly formed jaw bones would eventually lead to the ability to bite in some gnathostomes.
Newer research suggests that a branch of placoderms was most likely the ancestor of present-day gnathostomes. A 419-million-year-old fossil of a placoderm named Entelognathus had a bony skeleton and anatomical details associated with cartilaginous and bony fish, demonstrating that the absence of a bony skeleton in Chondrichthyes is a derived trait. The fossil findings of primitive bony fishes such as Guiyu oneiros and Psarolepis, which lived contemporaneously with Entelognathus and had pelvic girdles more in common with placoderms than with other bony fish, show that it was a relative rather than a direct ancestor of the extant gnathostomes. It also indicates that spiny sharks and Chondrichthyes represent a single sister group to the bony fishes. Fossil findings of juvenile placoderms, which had true teeth that grew on the surface of the jawbone and had no roots, making them impossible to replace or regrow as they broke or wore down as they grew older, proves the common ancestor of all gnathostomes had teeth and place the origin of teeth along with, or soon after, the evolution of jaws.
Late Ordovician-aged microfossils of what have been identified as scales of either acanthodians or "shark-like fishes", may mark Gnathostomata's first appearance in the fossil record. Undeniably unambiguous gnathostome fossils, mostly of primitive acanthodians, begin appearing by the early Silurian, and become abundant by the start of the Devonian.
Classification
Gnathostomata is traditionally an infraphylum, broken into three top-level groupings: Chondrichthyes, or the cartilaginous fish; Placodermi, an extinct grade of armored fish; and Teleostomi, which includes the familiar classes of bony fish, birds, mammals, reptiles, and amphibians. Some classification systems have used the term Amphirhina. It is a sister group of the jawless craniates Agnatha.
Evolution
The appearance of the early vertebrate jaw has been described as "a crucial innovation" and "perhaps the most profound and radical evolutionary step in the vertebrate history". Fish without jaws had more difficulty surviving than fish with jaws, and most jawless fish became extinct during the Triassic period. However studies of the cyclostomes, the jawless hagfishes and lampreys that did survive, have yielded little insight into the deep remodelling of the vertebrate skull that must have taken place as early jaws evolved.
The ancestor of all jawed vertebrates have gone through two rounds of whole genome duplication. The first happened before the gnathostome and cyclostome split, and appears to have been an autopolyploidy event (happened within the same species). The second occurred after the split, and was an allopolyploidy event (the result of hybridization between two lineages).
The customary view is that jaws are homologous to the gill arches. In jawless fishes a series of gills opened behind the mouth, and these gills became supported by cartilaginous elements. The first set of these elements surrounded the mouth to form the jaw. The upper portion of the second embryonic arch supporting the gill became the hyomandibular bone of jawed fish, which supports the skull and therefore links the jaw to the cranium. The hyomandibula is a set of bones found in the hyoid region in most fishes. It usually plays a role in suspending the jaws or the operculum in the case of teleosts.
While potentially older Ordovician records are known, the oldest unambigious evidence of jawed vertebrates are Qianodus and Fanjingshania from the early Silurian (Aeronian) of Guizhou, China around 439 million years ago, which are placed as acanthodian-grade stem-chondrichthyans.
| Biology and health sciences | Vertebrates, general classification | null |
211475 | https://en.wikipedia.org/wiki/Old%20World%20oriole | Old World oriole | The Old World orioles (Oriolidae) are an Old World family of passerine birds. The family contains 41 species which are divided in 4 genera. The family includes two extinct species from New Zealand that are placed in the genus Turnagra.
Taxonomy and systematics
The family Oriolidae comprises the piopios, figbirds, pitohuis and the Old World orioles. The piopios were added in 2011, having been formerly placed in the family Turnagridae. Several other genera have been proposed to split up the genus Oriolus. For example, the African black-headed species are sometimes placed in a separate genus, Baruffius. The family Oriolidae is not related to the New World orioles, despite their similar size, diet, behaviour and contrasting plumage patterns. Rather, these similarities are an example of convergent evolution.
Extant genera
There are three extant genera in the family Oriolidae:
Extinct genera
There are at least two extinct genera in the family Oriolidae:
Genus Turnagra – piopios (2 extinct species)
Genus Longmornis – Longmornis robustirostrata
Description
The orioles and figbirds are medium-sized passerines, around 20–30 cm in length, with the females only slightly smaller than the males. The beak is slightly curved and hooked, and, except in the figbirds, as long again as the head. The plumage of most species is bright and showy, although the females often have duller plumage than the males do. The plumage of many Australasian orioles mimics that of friarbirds (a genus of large honeyeaters), probably to reduce aggression against the smaller orioles.
Distribution and habitat
The family is distributed across Europe, Africa, Asia, and Australia. The few temperate nesting species are migratory, and some tropical species also show seasonal movements.
Behaviour and ecology
Breeding
Orioles are monogamous, breeding in territorial pairs (although the Australasian figbird, and possibly also the other figbirds, breed in loose colonies). Nesting sites may be chosen near aggressive species such as shrikes, drongos or friarbirds, which confer a degree of protection. The nest is a deep woven cup suspended like a hammock from a branch. They usually lay two or three eggs, but as many as six have been recorded.
Food and feeding
Orioles are arboreal and tend to feed in the canopy. Many species are able to survive in open forests and woodlands, although a few are restricted to closed forest. They are opportunistic omnivores, with the main components of their diet being fruit, berries, arthropods, and nectar.
| Biology and health sciences | Passerida | Animals |
211505 | https://en.wikipedia.org/wiki/Sepal | Sepal | A sepal () is a part of the flower of angiosperms (flowering plants). Usually green, sepals typically function as protection for the flower in bud, and often as support for the petals when in bloom.
Etymology
The term sepalum was coined by Noël Martin Joseph de Necker in 1790, and derived .
Collectively, the sepals are called the calyx (plural: calyces), the outermost whorl of parts that form a flower. The word calyx was adopted from the Latin , not to be confused with 'cup, goblet'. The Latin calyx is derived from Greek 'bud, calyx, husk, wrapping' ( Sanskrit 'bud'), while is derived from Greek 'cup, goblet'; both words have been used interchangeably in botanical Latin.
Description
The term tepal is usually applied when the parts of the perianth are difficult to distinguish, e.g. the petals and sepals share the same color or the petals are absent and the sepals are colorful. When the undifferentiated tepals resemble petals, they are referred to as "petaloid", as in petaloid monocots, orders of monocots with brightly colored tepals. Since they include Liliales, an alternative name is lilioid monocots. Examples of plants in which the term tepal is appropriate include genera such as Aloe and Tulipa. In contrast, genera such as Rosa and Phaseolus have well-distinguished sepals and petals.
The number of sepals in a flower is its merosity. Flower merosity is indicative of a plant's classification. The merosity of a eudicot flower is typically four or five. The merosity of a monocot or palaeodicot flower is three, or a multiple of three.
The development and form of the sepals vary considerably among flowering plants. They may be free (polysepalous) or fused together (gamosepalous). Often, the sepals are much reduced, appearing somewhat awn-like, or as scales, teeth, or ridges. Most often such structures protrude until the fruit is mature and falls off.
Examples of flowers with much-reduced perianths are found among the grasses.
In some flowers, the sepals are fused towards the base, forming a calyx tube (as in the family Lythraceae, and Fabaceae). In other flowers (e.g., Rosaceae, Myrtaceae), a hypanthium includes the bases of sepals, petals, and the attachment points of the stamens.
Mechanical cues may be responsible for sepal growth and there is a strong evidence suggesting that microtubules are present and determine the tensile strength and direction of growth at a molecular level.
Morphology
Morphologically, both sepals and petals are modified leaves. The calyx (the sepals) and the corolla (the petals) are the outer sterile whorls of the flower, which together form the perianth. In some plants, such as Aristolochia, the calyx is the primary whorl, forming a flower up to wide, with one sepal growing to a length of Aristolochia grandiflora, the largest of all calyces.
Function
Similarly to ordinary leaves, sepals are capable of performing photosynthesis. However, photosynthesis in sepals occurs at a slower rate than in ordinary leaves due to sepals having a lower stomatal density which limits the spaces for gas exchange.
After flowering, most plants have no more use for the calyx which withers or becomes vestigial, although in a few plants such as Lodoicea and eggplant (Solanum melongena) the calyx grows along with the fruit, possibly to protect the attachment point. Some plants retain a thorny calyx, either dried or live, as protection for the fruit or seeds. Examples include species of Acaena, some of the Solanaceae (for example the Tomatillo, Physalis philadelphica), and the water caltrop, Trapa natans. In some species, the calyx not only persists after flowering but instead of withering, begins to grow until it forms a bladder-like enclosure around the fruit. This is an effective protection against some kinds of birds and insects, for example in Hibiscus trionum and the Cape gooseberry. In other species, the calyx grows into an accessory fruit.
| Biology and health sciences | Plant anatomy and morphology: General | Biology |
211517 | https://en.wikipedia.org/wiki/Vulpecula | Vulpecula | Vulpecula is a faint constellation in the northern sky. Its name is Latin for "little fox", although it is commonly known simply as the fox. It was identified in the seventeenth century, and is located in the middle of the Summer Triangle (an asterism consisting of the bright stars Deneb, Vega, and Altair).
Features
Stars
There are no stars brighter than 4th magnitude in this constellation. The brightest star in Vulpecula is Alpha Vulpeculae, a magnitude 4.44m red giant at an approximate distance of 291 light-years. The star is an optical binary (separation of 413.7") that can be split using binoculars. The star also carries the traditional name Anser, which refers to the goose the little fox holds in its jaws.
23 Vulpeculae is the second brightest star in the constellation.
In 1967, the first pulsar, PSR B1919+21, was discovered in Vulpecula by Jocelyn Bell, supervised by Antony Hewish, in Cambridge. While they were searching for scintillation of radio signals of quasars, they observed pulses which repeated with a period of 1.3373 seconds. Terrestrial origin of the signal was ruled out because the time it took the object to reappear was a sidereal day instead of a solar day. This anomaly was finally identified as the signal of a rapidly rotating neutron star. Fifteen years after the first pulsar was discovered, the first millisecond pulsar, PSR B1937+21, was also discovered in Vulpecula, only a few degrees in the sky away from PSR B1919+21.
Vulpecula is also home to HD 189733 b, one of the closest extrasolar planets studied by the now-retired Spitzer Space Telescope. On 11 July 2007, a team led by Giovanna Tinetti published the results of their observations using the Spitzer Space Telescope concluding there is solid evidence for significant amounts of water vapor in the planet's atmosphere. Although HD 189733b with atmospheric temperatures rising above 1,000 °C is far from being habitable, this finding increases the likelihood that water, an essential component of life, would be found on a more Earth-like planet in the future.
Also located in Vulpecula is soft gamma repeater SGR 1935+2154. In 2020 it emitted a fast radio burst, the first one to be observed in the Milky Way.
Deep-sky objects
The Dumbbell Nebula (M27), is a large, bright planetary nebula which was discovered by the French astronomer Charles Messier in 1764 as the very first object of its kind. It can be seen with good binoculars in a dark sky location, appearing as a dimly glowing disk approximately 8 arcminutes in diameter. The nebula is approximately 9,800 years old. A telescope reveals its double-lobed shape, similar to that of an hourglass. Brocchi's Cluster (Collinder 399) is an asterism formerly thought to be an open cluster. It is also called "the Coathanger" because of its distinctive star pattern when viewed with binoculars or a low power telescope.
NGC 7052 is an elliptical galaxy in Vulpecula at a distance of 214 million light-years from Earth. It has a central dusty disk with a diameter of 3700 light-years; there is a supermassive black hole with a mass of 300 million solar masses in its nucleus. Astronomers surmise that the disk is the remnant of a smaller galaxy that merged with NGC 7052. Jets can be seen emanating from the galaxy, and it has very strong radio emissions. This means that it is also classified as a radio galaxy.
The eastern part of Vulpecula is occupied by the Hercules–Corona Borealis Great Wall. It is a galaxy filament, with the length of 3,000 megaparsecs, making it the largest known structure in the universe.
Stellar association
Vulpecula contains an OB-association of young stars, called the Vulpecula OB-association or Vul OB1. The association contains nearly 100 OB-stars and over 800 young stellar objects. It lies in the galactic plane, at a distance of about 2300 parsec. It contains the emission-line nebulae Sh-86, Sh-87 and Sh-88. Vul OB1 has a length of about 100 parsec and is sculpting many pillar-like structures in this region.
History
In the late 17th century, the astronomer Johannes Hevelius created Vulpecula. It was originally known as Vulpecula cum ansere ("the little fox with the goose") or Vulpecula et Anser ("the little fox and the goose"), and was illustrated with a goose in the jaws of a fox. Hevelius did not regard the fox and the goose to be two separate constellations, but later the stars were divided into a separate Anser and Vulpecula. Today, they have been merged again under the name of the fox, but the goose is remembered by the name of the star α Vulpeculae: Anser.
| Physical sciences | Other | Astronomy |
211602 | https://en.wikipedia.org/wiki/Hidden-variable%20theory | Hidden-variable theory | In physics, a hidden-variable theory is a deterministic physical model which seeks to explain the probabilistic nature of quantum mechanics by introducing additional (possibly inaccessible) variables.
Indeterminacy of the state of a system previous to measurement is assumed to be a part of the mathematical formulation of quantum mechanics; moreover, bounds for indeterminacy can be expressed in a quantitative form by the Heisenberg uncertainty principle. Most hidden-variable theories are attempts to avoid this indeterminacy, but possibly at the expense of requiring that nonlocal interactions be allowed. One notable hidden-variable theory is the de Broglie–Bohm theory.
In 1935, Albert Einstein, Boris Podolsky, and Nathan Rosen in their EPR paper argued that quantum entanglement might indicate quantum mechanics is an incomplete description of reality. John Stewart Bell in 1964, in his eponymous theorem proved that correlations between particles under any local hidden variable theory must obey certain constraints. Subsequently, Bell test experiments have demonstrated broad violation of these constraints, ruling out such theories. Bell's theorem, however, does not rule out the possibility of nonlocal theories or superdeterminism; these therefore cannot be falsified by Bell tests.
Motivation
Macroscopic physics requires classical mechanics which allows accurate predictions of mechanical motion with reproducible, high precision. Quantum phenomena require quantum mechanics, which allows accurate predictions of statistical averages only. If quantum states had hidden-variables awaiting ingenious new measurement technologies, then the latter (statistical results) might be convertible to a form of the former (classical-mechanical motion).
Such a classical mechanics would eliminate unsettling characteristics of quantum theory like the uncertainty principle. More fundamentally however, a successful model of quantum phenomena with hidden variables implies quantum entities with intrinsic values independent of measurements. Existing quantum mechanics asserts that state properties can only be known after a measurement. As N. David Mermin puts it:
In other words, whereas a hidden-variable theory would imply intrinsic particle properties, in quantum mechanics an electron has no definite position and velocity to even be revealed.
History
"God does not play dice"
In June 1926, Max Born published a paper, in which he was the first to clearly enunciate the probabilistic interpretation of the quantum wave function, which had been introduced by Erwin Schrödinger earlier in the year. Born concluded the paper as follows:
Here the whole problem of determinism comes up. From the standpoint of our quantum mechanics there is no quantity which in any individual case causally fixes the consequence of the collision; but also experimentally we have so far no reason to believe that there are some inner properties of the atom which conditions a definite outcome for the collision. Ought we to hope later to discover such properties ... and determine them in individual cases? Or ought we to believe that the agreement of theory and experiment—as to the impossibility of prescribing conditions for a causal evolution—is a pre-established harmony founded on the nonexistence of such conditions? I myself am inclined to give up determinism in the world of atoms. But that is a philosophical question for which physical arguments alone are not decisive.
Born's interpretation of the wave function was criticized by Schrödinger, who had previously attempted to interpret it in real physical terms, but Albert Einstein's response became one of the earliest and most famous assertions that quantum mechanics is incomplete:
Quantum mechanics is very worthy of respect. But an inner voice tells me this is not the genuine article after all. The theory delivers much but it hardly brings us closer to the Old One's secret. In any event, I am convinced that He is not playing dice.
Niels Bohr reportedly replied to Einstein's later expression of this sentiment by advising him to "stop telling God what to do."
Early attempts at hidden-variable theories
Shortly after making his famous "God does not play dice" comment, Einstein attempted to formulate a deterministic counter proposal to quantum mechanics, presenting a paper at a meeting of the Academy of Sciences in Berlin, on 5 May 1927, titled "Bestimmt Schrödinger's Wellenmechanik die Bewegung eines Systems vollständig oder nur im Sinne der Statistik?" ("Does Schrödinger's wave mechanics determine the motion of a system completely or only in the statistical sense?"). However, as the paper was being prepared for publication in the academy's journal, Einstein decided to withdraw it, possibly because he discovered that, contrary to his intention, his use of Schrödinger's field to guide localized particles allowed just the kind of non-local influences he intended to avoid.
At the Fifth Solvay Congress, held in Belgium in October 1927 and attended by all the major theoretical physicists of the era, Louis de Broglie presented his own version of a deterministic hidden-variable theory, apparently unaware of Einstein's aborted attempt earlier in the year. In his theory, every particle had an associated, hidden "pilot wave" which served to guide its trajectory through space. The theory was subject to criticism at the Congress, particularly by Wolfgang Pauli, which de Broglie did not adequately answer; de Broglie abandoned the theory shortly thereafter.
Declaration of completeness of quantum mechanics, and the Bohr–Einstein debates
Also at the Fifth Solvay Congress, Max Born and Werner Heisenberg made a presentation summarizing the recent tremendous theoretical development of quantum mechanics. At the conclusion of the presentation, they declared:
[W]hile we consider ... a quantum mechanical treatment of the electromagnetic field ... as not yet finished, we consider quantum mechanics to be a closed theory, whose fundamental physical and mathematical assumptions are no longer susceptible of any modification....
On the question of the 'validity of the law of causality' we have this opinion: as long as one takes into account only experiments that lie in the domain of our currently acquired physical and quantum mechanical experience, the assumption of indeterminism in principle, here taken as fundamental, agrees with experience.
Although there is no record of Einstein responding to Born and Heisenberg during the technical sessions of the Fifth Solvay Congress, he did challenge the completeness of quantum mechanics at various times. In his tribute article for Born's retirement he discussed the quantum representation of a macroscopic ball bouncing elastically between rigid barriers. He argues that such a quantum representation does not represent a specific ball, but "time ensemble of systems". As such the representation is correct, but incomplete because it does not represent the real individual macroscopic case. Einstein considered quantum mechanics incomplete "because the state function, in general, does not even describe the individual event/system".
Von Neumann's proof
John von Neumann in his 1932 book Mathematical Foundations of Quantum Mechanics had presented a proof that there could be no "hidden parameters" in quantum mechanics. The validity of von Neumann's proof was questioned by Grete Hermann in 1935, who found a flaw in the proof. The critical issue concerned averages over ensembles. Von Neumann assumed that a relation between the expected values of different observable quantities holds for each possible value of the "hidden parameters", rather than only for a statistical average over them. However Hermann's work went mostly unnoticed until its rediscovery by John Stewart Bell more than 30 years later.
The validity and definitiveness of von Neumann's proof were also questioned by Hans Reichenbach, and possibly in conversation though not in print by Albert Einstein. Reportedly, in a conversation circa 1938 with his assistants Peter Bergmann and Valentine Bargmann, Einstein pulled von Neumann's book off his shelf, pointed to the same assumption critiqued by Hermann and Bell, and asked why one should believe in it. Simon Kochen and Ernst Specker rejected von Neumann's key assumption as early as 1961, but did not publish a criticism of it until 1967.
EPR paradox
Einstein argued that quantum mechanics could not be a complete theory of physical reality. He wrote,
Consider a mechanical system consisting of two partial systems A and B which interact with each other only during a limited time. Let the ψ function [i.e., wavefunction] before their interaction be given. Then the Schrödinger equation will furnish the ψ function after the interaction has taken place. Let us now determine the physical state of the partial system A as completely as possible by measurements. Then quantum mechanics allows us to determine the ψ function of the partial system B from the measurements made, and from the ψ function of the total system. This determination, however, gives a result which depends upon which of the physical quantities (observables) of A have been measured (for instance, coordinates or momenta). Since there can be only one physical state of B after the interaction which cannot reasonably be considered to depend on the particular measurement we perform on the system A separated from B it may be concluded that the ψ function is not unambiguously coordinated to the physical state. This coordination of several ψ functions to the same physical state of system B shows again that the ψ function cannot be interpreted as a (complete) description of a physical state of a single system.
Together with Boris Podolsky and Nathan Rosen, Einstein published a paper that gave a related but distinct argument against the completeness of quantum mechanics. They proposed a thought experiment involving a pair of particles prepared in what would later become known as an entangled state. Einstein, Podolsky, and Rosen pointed out that, in this state, if the position of the first particle were measured, the result of measuring the position of the second particle could be predicted. If instead the momentum of the first particle were measured, then the result of measuring the momentum of the second particle could be predicted. They argued that no action taken on the first particle could instantaneously affect the other, since this would involve information being transmitted faster than light, which is impossible according to the theory of relativity. They invoked a principle, later known as the "EPR criterion of reality", positing that: "If, without in any way disturbing a system, we can predict with certainty (i.e., with probability equal to unity) the value of a physical quantity, then there exists an element of reality corresponding to that quantity." From this, they inferred that the second particle must have a definite value of both position and of momentum prior to either quantity being measured. But quantum mechanics considers these two observables incompatible and thus does not associate simultaneous values for both to any system. Einstein, Podolsky, and Rosen therefore concluded that quantum theory does not provide a complete description of reality.
Bohr answered the Einstein–Podolsky–Rosen challenge as follows:
[The argument of] Einstein, Podolsky and Rosen contains an ambiguity as regards the meaning of the expression "without in any way disturbing a system." ... [E]ven at this stage [i.e., the measurement of, for example, a particle that is part of an entangled pair], there is essentially the question of an influence on the very conditions which define the possible types of predictions regarding the future behavior of the system. Since these conditions constitute an inherent element of the description of any phenomenon to which the term "physical reality" can be properly attached, we see that the argumentation of the mentioned authors does not justify their conclusion that quantum-mechanical description is essentially incomplete."
Bohr is here choosing to define a "physical reality" as limited to a phenomenon that is immediately observable by an arbitrarily chosen and explicitly specified technique, using his own special definition of the term 'phenomenon'. He wrote in 1948:
As a more appropriate way of expression, one may strongly advocate limitation of the use of the word phenomenon to refer exclusively to observations obtained under specified circumstances, including an account of the whole experiment.Rosenfeld, L. (). 'Niels Bohr's contribution to epistemology', pp. 522–535 in Selected Papers of Léon Rosenfeld, Cohen, R.S., Stachel, J.J. (editors), D. Riedel, Dordrecht, , p. 531: "Moreover, the complete definition of the phenomenon must essentially contain the indication of some permanent mark left upon a recording device which is part of the apparatus; only by thus envisaging the phenomenon as a closed event, terminated by a permanent record, can we do justice to the typical wholeness of the quantal processes."
This was, of course, in conflict with the EPR criterion of reality.
Bell's theorem
In 1964, John Stewart Bell showed through his famous theorem that if local hidden variables exist, certain experiments could be performed involving quantum entanglement where the result would satisfy a Bell inequality. If, on the other hand, statistical correlations resulting from quantum entanglement could not be explained by local hidden variables, the Bell inequality would be violated. Another no-go theorem concerning hidden-variable theories is the Kochen–Specker theorem.
Physicists such as Alain Aspect and Paul Kwiat have performed experiments that have found violations of these inequalities up to 242 standard deviations. This rules out local hidden-variable theories, but does not rule out non-local ones. Theoretically, there could be experimental problems that affect the validity of the experimental findings.
Gerard 't Hooft has disputed the validity of Bell's theorem on the basis of the superdeterminism loophole and proposed some ideas to construct local deterministic models.
Bohm's hidden-variable theory
In 1952, David Bohm proposed a hidden variable theory. Bohm unknowingly rediscovered (and extended) the idea that Louis de Broglie's pilot wave theory had proposed in 1927 (and abandoned) – hence this theory is commonly called "de Broglie-Bohm theory". Assuming the validity of Bell's theorem, any deterministic hidden-variable theory that is consistent with quantum mechanics would have to be non-local, maintaining the existence of instantaneous or faster-than-light relations (correlations) between physically separated entities.
Bohm posited both the quantum particle, e.g. an electron, and a hidden 'guiding wave' that governs its motion. Thus, in this theory electrons are quite clearly particles. When a double-slit experiment is performed, the electron goes through either one of the slits. Also, the slit passed through is not random but is governed by the (hidden) pilot wave, resulting in the wave pattern that is observed.
In Bohm's interpretation, the (non-local) quantum potential constitutes an implicate (hidden) order which organizes a particle, and which may itself be the result of yet a further implicate order: a superimplicate order which organizes a field. Nowadays Bohm's theory is considered to be one of many interpretations of quantum mechanics. Some consider it the simplest theory to explain quantum phenomena. Nevertheless, it is a hidden-variable theory, and necessarily so. The major reference for Bohm's theory today is his book with Basil Hiley, published posthumously.
A possible weakness of Bohm's theory is that some (including Einstein, Pauli, and Heisenberg) feel that it looks contrived. (Indeed, Bohm thought this of his original formulation of the theory.) Bohm said he considered his theory to be unacceptable as a physical theory due to the guiding wave's existence in an abstract multi-dimensional configuration space, rather than three-dimensional space.
Recent developments
In August 2011, Roger Colbeck and Renato Renner published a proof that any extension of quantum mechanical theory, whether using hidden variables or otherwise, cannot provide a more accurate prediction of outcomes, assuming that observers can freely choose the measurement settings. Colbeck and Renner write: "In the present work, we have ... excluded the possibility that any extension of quantum theory (not necessarily in the form of local hidden variables) can help predict the outcomes of any measurement on any quantum state. In this sense, we show the following: under the assumption that measurement settings can be chosen freely, quantum theory really is complete".
In January 2013, Giancarlo Ghirardi and Raffaele Romano described a model which, "under a different free choice assumption [...] violates [the statement by Colbeck and Renner] for almost all states of a bipartite two-level system, in a possibly experimentally testable way".
| Physical sciences | Quantum mechanics | Physics |
211647 | https://en.wikipedia.org/wiki/Giardia%20duodenalis | Giardia duodenalis | Giardia duodenalis, also known as Giardia intestinalis and Giardia lamblia, is a flagellated parasitic protozoan microorganism of the genus Giardia that colonizes the small intestine, causing a diarrheal condition known as giardiasis. The parasite attaches to the intestinal epithelium by a ventral disc (syn. adhesive disc or sucker), and reproduces via binary fission. G. duodenalis is a non-invasive parasite, that does not spread to other parts of the gastrointestinal tract, but remains confined to the lumen of the small intestine. The parasite exists in two forms; trophozoites and cysts. The microorganism can undergo encystation, transforming into a dormant cyst that enables it to survive outside of its host. Giardia trophozoites are anaerobic, and absorb their nutrients from the intestinal lumen. If the organism is stained, its characteristic pattern resembles the familiar "smiley face" symbol.
Chief pathways of human infection include ingestion of untreated drinking water (which is the most common method of transmission for this parasite), food, soil contaminated with human feces, and sewage, a phenomenon particularly common in many developing countries. Contamination of natural waters also occurs in watersheds where intensive grazing occurs.
Giardia infections occur worldwide. It is the most commonly identified intestinal parasite among children in day-care centers, hikers and immunocompromised patients. About 20,000 cases per year in the United States are reported.
Almost half of those infected with giardiasis remain asymptomatic. For those who do experience symptoms, they usually appear 1 to 2 weeks after infection. Common symptoms include abdominal pain, nausea, and bloating, along with large, watery, foul-smelling, and greasy stools. Due to frequent loose stools, individuals with giardiasis often experience dehydration. It has also been shown that G. intestinalis damages the intestinal epithelium, which directly affects nutrient absorption. In severe cases, giardiasis can lead to chronic diarrhea, chronic fatigue syndrome and cognitive impairment in children.
Life cycle
G. duodenalis takes on two morphologically distinct forms during its lifecycle. Trophozoites are the replicative stage of the parasite, characterized by a pear-shaped, motile, flagellated cell that survives only in the small intestine of the host. The trophozoites do not penetrate host cells, but rather attaches to the intestinal epithelium cells to establish an infection. A cyst is the environmentally stable stage of the parasite, that facilitates transmission between hosts.
The infection process in the host begins with the ingestion of cysts, which pass through the stomach into the first part of small intestine, or the duodenum. Exposure to digestive enzymes and an acidic pH triggers excystation, where the cysts release trophozoites. Trophozoites swim through the intestinal mucus until they eventually adhere to the intestinal epithelium using their ventral disc. Adhered trophozoites can then feed, divide by binary fission and cause the disease. Trophozoites cause structural and functional damage to the host epithelial cells, impairing the intestine's ability to absorb nutrients effectively. Some trophozoites differentiate back into cysts under specific conditions, such as high organism density. As the trophozoites travel through the intestinal lumen to the large intestine, the alkaline environment and the presence of pancreatic proteases trigger the encystation. Trophozoites develop into cysts through encystation, which involves specialized vesicles (ESVs) that facilitate the formation of the cyst wall.
Cysts and trophozoites pass through the host's large intestine and are shed in the feces. Trophozoites cannot survive outside of the host, whereas cysts can remain viable in the external environment for several months. The cysts remain dormant until ingested by a host animal. When a new potential host ingests water or food contaminated with this feces, the cysts gain entry to the gastrointestinal tract of the new host, repeating the cycle.
Structure
The trophozoite has an elaborate structure with two nuclei and four pairs of flagella which allow it to swim within the intestinal lumen of the host. It also has an adhesive disk on its ventral surface that is associated with the parasite's attachment to the intestinal epithelium. The adhesive disk is composed of microtubules, and is found only in Giardia. The parasite lacks Golgi apparatus and mitochondria but has mitosomes, which probably evolved from mitochondria. The mitosome is a double-membraned organelle, that lacks the enzymatic components required for classic mitochondrial functions, such as ATP synthesis and lipid metabolism. However, they do contain certain mitochondrial genes associated with iron-sulfur complex biosynthesis, which suggests that this organism likely lost its mitochondria during evolution.
The cyst contains four nuclei, axonemes, median bodies as well as the rigid cyst wall. Additionally, the cysts contain fragments of the adhesive disc within the cytoplasm. The two-layered cyst wall protects the parasite against different environmental conditions such as chemical treatments. The metabolic activity of the cyst is significantly lower compared to that of the trophozoites, which allows the parasite to survive for longer periods in harsh conditions, such as in cold environments.
Geographical Prevalence of Giardia duodenalis
G. duodenalis is the most widespread intestinal parasite affecting humans. The parasite Giardia duodenalis can be found all over the world, in both developing and industrialized nations. However, it is most commonly found in tropical and temperate climates. Giardia duodenalis is common around the world because the parasite resides in bodies of water; typically rivers, lakes, and recreational swimming pools. Giardiasis is more prevalent in developing countries, where the sanitation and overall hygiene is poorer compared to developed countries. In developed nations, giardiasis has a prevalence of 2%-5%, whereas in developing nations it is significantly higher, ranging from 20% to 30%. In the United States, it has been discovered that a majority of whom are infected by the Giardia duodenalis parasite tend to reside in more urban areas, and, patients who are infected are more likely to live in the Southern United States.
Prevalence and Epidemiology
G. duodenalis causes an infection called giardiasis. This disease is the cause of both endemic and epidemic disease worldwide and is the most frequently identified intestinal parasite in the United States and Canada. An infected individual can excrete between 1 million and 1 billion cysts daily, and the infectious dose can be as low as 10 cysts. This makes Giardia extremely infectious.
It is estimated to infect over 280 million people world every yearresulting over 500,000 deaths. The most affected demographic is children 0 to 4 years of age. Globally G. duodenalis is the most commonly identified protozoal intestinal parasite. Giardia has common seasonal patterns in the distribution of infection rates with highest peaks in the late summer to early fall.
The cyst can survive for weeks to months in cold water, so the parasite can be present in contaminated wells and water systems, especially in stagnant water sources, such as naturally occurring ponds, storm-water storage systems, and even clean-looking mountain streams. Cysts can also be found on contaminated surfaces, soil and food. They may also occur in city reservoirs and persist after water treatment, as the cysts are resistant to conventional water-treatment methods, such as chlorination and ozonolysis. Zoonotic transmission is also possible, but is less frequent. Giardia infection is a concern for people camping in the wilderness or swimming in contaminated streams or lakes, especially the artificial lakes formed by beaver dams (hence the popular name for giardiasis, "beaver fever").
In addition to waterborne sources, Giardia infections are more commonly found in children compared to adults, this is believed to be due to fecal-oral transmission of the cysts. For example, in developed countries it affects approximately 2% of adults and 8% of children. In developing countries the prevalence rates reach 15% to 20% in children under 10 years old. Thus, there is a significant variation in infection rates based on geographical area.
Those who work with children are also at risk of being infected, as are family members of infected individuals. 7% of children aged 1 to 3 years and 11% of infants and toddlers tested for admission to day-care centers were found to be infected. Not all Giardia infections are symptomatic, and many people can unknowingly serve as carriers of the parasite. Re- infection and chronic infections of the parasite can occur.
Ecology
Giardia infects humans, but is also one of the most common parasites infecting cats, dogs, and birds. Mammalian hosts also include dozens of species, including cattle, sheep, and goats.
Cats can be cured easily, and lambs usually simply lose weight, but in calves, the parasites can be fatal and often are not responsive to antibiotics or electrolytes. Carriers among calves can also be asymptomatic. This parasite is deadly for chinchillas, so extra care must be taken by providing them with safe water. Dogs have a high infection rate, as 30% of the population under one year old are known to be infected in kennels. The infection is more prevalent in puppies than in adult dogs. Infected dogs can be isolated and treated, or the entire pack at a kennel can be presumptively treated together. Kennels and areas used for exercise should be considered contaminated for at least one month after dogs show signs of infection, as cysts can survive in the environment for long periods of time. Prevention can be achieved by quarantine of infected dogs for at least 20 days and careful management and maintenance of a clean water supply.
Cell biology
G. duodenalis trophozoites are pear-shaped cells, 10 to 20 μm long, 7 to 10 μm across, and 2 to 4 μm thick. They are motile by way of four pairs of flagella, which propel the trophozoites through the intestine. Notably, each G. duodenalis cell has two nuclei, both of which actively transcribe genes. Adjacent to the nucleus, G. duodenalis cells have an endoplasmic reticulum that extends through much of the cell. Trophozoites about to differentiate into cysts also contain prominent vesicles termed encystation-specific vesicles that disappear once cyst wall construction begins. Unlike most other eukaryotes, G. duodenalis cells contain no visible mitochondria, but instead contains a substantially reduced metabolic organelle known as a mitosome. Additionally, cells appear to contain no Golgi bodies, and instead the secretory system consists entirely of the endoplasmic reticulum and numerous vesicles dispersed throughout the cell, termed peripheral vesicles. Peripheral vesicles are responsible both for taking up extracellular nutrients, and expelling waste outside the cell. Each cell also contains a pair of rigid structures called median bodies which make up part of the G. lamblia cytoskeleton. Trophozoites adhere to host epithelial cells via a specialized disk-shaped organelle called the ventral disk.
Cysts are oval-shaped cells slightly smaller than trophozoites. They lack flagella, and are covered by a smooth, clear cyst wall. Each cyst contains four nuclei and fragments of the ventral disc.
Metabolism
G. duodenalis primarily generates its energy by breaking down glucose via glycolysis, as well as the arginine deiminase pathway. It is unable to synthesize nucleotides on its own, instead salvaging them from its host. Synthesis of iron–sulfur clusters is done in a double-membrane-bound compartment called the mitosome, which is likely a remnant of mitochondria. Each cell contains 25 to 100 mitosomes divided into two categories - peripheral mitosomes, which are scattered throughout the cell, and central mitosomes, which gather at the center of the cell for unknown reasons. As in mitochondria, proteins with a certain peptide signal sequence are trafficked to and imported into the mitosome. Unlike mitochondria, mitosomes have no genome of their own. All mitosomal genes are encoded by the Giardia nuclear genome.
Genetics
Giardia and the other diplomonads are unique in their possession of two cell nuclei that are similar in appearance, DNA content, transcription, and time of replication. Giardia is a polyploid organism, with at least four, and perhaps eight or more, copies of each of five chromosomes per organism. The genome has been sequenced and was published in 2007, although the sequence contains several gaps. The sequence is about 12 million base pairs and contains about 5000 protein-coding genes. The GC-content is 46%. Trophozoites have a ploidy of four and the ploidy of cysts is eight, which in turn raises the question of how Giardia maintains homogeneity between the chromosomes of the same and opposite nuclei. Modern sequencing technologies have been used to resequence different strains.
Immunology
Infections with Giardia are self-limited in immunocompetent individuals, while people with immunodeficiency disorders may develop chronic giardiasis. During the infection different mechanisms from the innate and adaptive immune system are activated. The first physical barrier is the mucus layer where the organism interacts with epithelial, immune cells, and some antimicrobial peptides released by those cells as well as nitric oxide and inflammatory cytokines like interleukin 6. TLR2 and TLR4 also can be activated by Giardia. The T-cell response in giardiasis includes T helper cells and cytotoxic T cells, and the production of IgA by B cells also helps to eliminate the infection.
Evolution
Giardia had been assumed to be primitively asexual and with no means of transferring DNA between nuclei. These assumptions made explaining the remarkably low level of allelic heterozygosity (< 0.01%) in the genome isolate, WB, very difficult, but all those assumptions of asexuality are now in doubt, with population genetics providing evidence for recombination and the identification of meiotic genes, evidence for recombination among isolates and the evidence for exchange of genetic material between nuclei during the process of encystation.
These findings on sexuality in Giardia, above, have important implications for understanding the origin of sexual reproduction in eukaryotes. Though sexual reproduction is widespread among extant eukaryotes, until recently, sex seemed unlikely to be a primordial and fundamental feature of eukaryotes. A probable reason for the view that sex may not be fundamental to eukaryotes was that sexual reproduction previously appeared to be lacking in certain human pathogenic single-celled eukaryotes (e.g. Giardia) that diverged from early ancestors in the eukaryotic lineage.
In addition to the evidence cited above for recombination in Giardia, Malik et al. reported that many meiosis specific genes occur in the Giardia genome, and further that homologs of these genes also occur in another unicellular eukaryote, Trichomonas vaginalis. Because these two species are descendants of lineages that are highly divergent among eukaryotes, Malik et al. suggested that these meiotic genes were present in a common ancestor of all eukaryotes. Thus, on this view, the earliest ancestor of eukaryotes was likely capable of sexual reproduction. Furthermore, Dacks and Roger proposed, based on phylogenetic analysis, that facultative sex was present in the common ancestor of all eukaryotes. Bernstein et al. also reviewed evidence in support of this view.
Eight genotype assemblages of G. duodenalis have been recognized to date (A-H). Genotyping of G. duodenalis isolated from various hosts has shown that assemblages A and B infect the largest range of host species, and appear to be the main (or possibly only) G. duodenalis assemblages that undeniably infect human subjects.
Research
Frances Gillin of the University of California, San Diego, and her colleagues cultivated the entire lifecycle of this parasite in the laboratory, and identified biochemical cues in the host's digestive system that trigger Giardia'''s lifecycle transformations. They also uncovered several ways in which the parasite evades the defenses of the infected organism. One of these is by altering the proteins on its surface, which confounds the ability of the infected animal's immune system to detect and combat the parasite (called antigenic variation). Gillin's work reveals why Giardia infections are extremely persistent and prone to recur. In addition, these insights into its biology and survival techniques may enable scientists to develop better strategies to understand, prevent, and treat Giardia infections.
In December 2008, Nature published an article showing the discovery of an RNA interference mechanism that allows Giardia to switch variant-specific surface proteins to avoid host immune response. The discovery was made by the team working at the Biochemistry and Molecular Biology Laboratory, School of Medicine, Catholic University of Cordoba, Argentina, led by Dr. Hugo Lujan.
The main congress about Giardia is the International Giardia and Cryptosporidium Conference. A summary of results presented at the most recent edition (2019, in Rouen, France) is available.
In 2022, a study conducted by Elisa Barroeta-Echegaray and colleagues concluded that Giardia duodenalis secretes enolase as a monomer during the interaction, or attachment, of trophozoites with intestinal epithelial cells. This interaction was shown to activate plasminogen and induce necroptotic damage in intestinal epithelial cells. The researchers demonstrated that blocking the enolase inhibited trophozoite attachment to intestinal epithelial cells. Furthermore, enolase was shown to enhance plasmin activity, leading to significant cell damage characterized by vacuolization and intercellular separation. Enolase also induced necroptosis in epithelial cells via tumor necrosis factor α (TNF-α) and apoptosis-inducing factor (AIF), independent of caspase-3 activity. These findings suggest that Giardia enolase is a critical virulence factor in host-pathogen interactions.
History
The first likely description of Giardia was in 1681 by Antonie van Leeuwenhoek, who in a letter to Robert Hooke, described "animalcules" resembling Giardia trophozoites in his stool. The next known description of Giardia wasn't until 1859, when Czech physician Vilém Lambl published a description of the trophozoite stages he saw in the stool of a pediatric patient. Lambl termed the organism Cercomonas intestinalis. In 1888, Raphaël Blanchard renamed the parasite Lamblia intestinalis in Lambl's honor. In 1915, Charles Stiles renamed the organism Giardia lamblia'' in honor of both Lambl and Professor Alfred Mathieu Giard of Paris. In 1921, Charles E. Simon published a detailed description of the parasite's morphology.
| Biology and health sciences | Excavata | Plants |
211724 | https://en.wikipedia.org/wiki/Epoxy | Epoxy | Epoxy is the family of basic components or cured end products of epoxy resins. Epoxy resins, also known as polyepoxides, are a class of reactive prepolymers and polymers which contain epoxide groups. The epoxide functional group is also collectively called epoxy. The IUPAC name for an epoxide group is an oxirane.
Epoxy resins may be reacted (cross-linked) either with themselves through catalytic homopolymerisation, or with a wide range of co-reactants including polyfunctional amines, acids (and acid anhydrides), phenols, alcohols and thiols (sometimes called mercaptans). These co-reactants are often referred to as hardeners or curatives, and the cross-linking reaction is commonly referred to as curing.
Reaction of polyepoxides with themselves or with polyfunctional hardeners forms a thermosetting polymer, often with favorable mechanical properties and high thermal and chemical resistance. Epoxy has a wide range of applications, including metal coatings, composites, use in electronics, electrical components (e.g. for chips on board), LEDs, high-tension electrical insulators, paintbrush manufacturing, fiber-reinforced plastic materials, and adhesives for structural and other purposes.
The health risks associated with exposure to epoxy resin compounds include contact dermatitis and allergic reactions, as well as respiratory problems from breathing vapor and sanding dust, especially from compounds not fully cured.
History
Condensation of epoxides and amines was first reported and patented by Paul Schlack of Germany in 1934. Claims of discovery of bisphenol-A-based epoxy resins include Pierre Castan in 1943. Castan's work was licensed by Ciba, Ltd. of Switzerland, which went on to become one of the three major epoxy resin producers worldwide. In 1946, Sylvan Greenlee, working for the Devoe & Raynolds Company (now part of Hexion Inc.), patented resin derived from bisphenol-A and epichlorohydrin.
Chemistry
Most of the commercially used epoxy monomers are produced by the reaction of a compound with acidic hydroxy groups and epichlorohydrin. First a hydroxy group reacts in a coupling reaction with epichlorohydrin, followed by dehydrohalogenation. Epoxy resins produced from such epoxy monomers are called glycidyl-based epoxy resins. The hydroxy group may be derived from aliphatic diols, polyols (polyether polyols), phenolic compounds or dicarboxylic acids. Phenols can be compounds such as bisphenol A and novolak. Polyols can be compounds such as 1,4-butanediol. Di- and polyols lead to glycidyl ethers. Dicarboxylic acids such as hexahydrophthalic acid are used for diglycide ester resins. Instead of a hydroxy group, also the nitrogen atom of an amine or amide can be reacted with epichlorohydrin.
The other production route for epoxy resins is the conversion of aliphatic or cycloaliphatic alkenes with peracids: In contrast to glycidyl-based epoxy resins, this production of such epoxy monomers does not require an acidic hydrogen atom but an aliphatic double bond.
The epoxide group is also sometimes referred to as an oxirane group.
Bisphenol-based
The most common epoxy resins are based on reacting epichlorohydrin (ECH) with bisphenol A, resulting in a different chemical substance known as bisphenol A diglycidyl ether (commonly known as BADGE or DGEBA). Bisphenol A-based resins are the most widely commercialised resins but also other bisphenols are analogously reacted with epichlorohydrin, for example Bisphenol F.
In this two-stage reaction, epichlorohydrin is first added to bisphenol A (bis(3-chloro-2-hydroxy-propoxy)bisphenol A is formed), then a bisepoxide is formed in a condensation reaction with a stoichiometric amount of sodium hydroxide. The chlorine atom is released as sodium chloride (NaCl) and the hydrogen atom as water.
Higher molecular weight diglycidyl ethers (n ≥ 1) are formed by the reaction of the bisphenol A diglycidyl ether formed with further bisphenol A, this is called prepolymerization:
A product comprising a few repeat units (n = 1 to 2) is a viscous, clear liquid; this is called a liquid epoxy resin. A product comprising more repeating units (n = 2 to 30) is at room temperature a colourless solid, which is correspondingly referred to as solid epoxy resin.
Instead of bisphenol A, other bisphenols (especially bisphenol F) or brominated bisphenols (e. g. tetrabromobisphenol A) can be used for the said epoxidation and prepolymerisation. Bisphenol F may undergo epoxy resin formation in a similar fashion to bisphenol A. These resins typically have lower viscosity and a higher mean epoxy content per gram than bisphenol A resins, which (once cured) gives them increased chemical resistance.
Important epoxy resins are produced from combining epichlorohydrin and bisphenol A to give bisphenol A diglycidyl ethers.
Increasing the ratio of bisphenol A to epichlorohydrin during manufacture produces higher molecular weight linear polyethers with glycidyl end groups, which are semi-solid to hard crystalline materials at room temperature depending on the molecular weight achieved. This route of synthesis is known as the "taffy" process. The usual route to higher molecular weight epoxy resins is to start with liquid epoxy resin (LER) and add a calculated amount of bisphenol A and then a catalyst is added and the reaction heated to circa . This process is known as "advancement". As the molecular weight of the resin increases, the epoxide content reduces and the material behaves more and more like a thermoplastic. Very high molecular weight polycondensates (ca. 30,000–70,000 g/mol) form a class known as phenoxy resins and contain virtually no epoxide groups (since the terminal epoxy groups are insignificant compared to the total size of the molecule). These resins do however contain hydroxyl groups throughout the backbone, which may also undergo other cross-linking reactions, e.g. with aminoplasts, phenoplasts and isocyanates.
Epoxy resins are polymeric or semi-polymeric materials or an oligomer, and as such rarely exist as pure substances, since variable chain length results from the polymerisation reaction used to produce them. High purity grades can be produced for certain applications, e.g. using a distillation purification process. One downside of high purity liquid grades is their tendency to form crystalline solids due to their highly regular structure, which then require melting to enable processing.
An important criterion for epoxy resins is the Epoxy value which is connected to the epoxide group content. This is expressed as the "epoxide equivalent weight", which is the ratio between the molecular weight of the monomer and the number of epoxide groups. This parameter is used to calculate the mass of co-reactant (hardener) to use when curing epoxy resins. Epoxies are typically cured with stoichiometric or near-stoichiometric quantities of hardener to achieve the best physical properties.
Novolaks
Novolaks are produced by reacting phenol with methanal (formaldehyde). The reaction of epichlorohydrin and novolaks produces novolaks with glycidyl residues, such as epoxyphenol novolak (EPN) or epoxycresol novolak (ECN). These highly viscous to solid resins typically carry 2 to 6 epoxy groups per molecule. By curing, highly cross-linked polymers with high temperature and chemical resistance but low mechanical flexibility are formed due to the high functionality, and hence high crosslink density of these resins.
Aliphatic
There are two common types of aliphatic epoxy resins: those obtained by epoxidation of double bonds (cycloaliphatic epoxides and epoxidized vegetable oils) and those formed by reaction with epichlorohydrin (glycidyl ethers and esters).
Cycloaliphatic epoxides contain one or more aliphatic rings in the molecule on which the oxirane ring is contained (e.g. 3,4-epoxycyclohexylmethyl-3',4'-epoxycyclohexane carboxylate). They are produced by the reaction of a cyclic alkene with a peracid (see above). Cycloaliphatic epoxides are characterised by their aliphatic structure, high oxirane content and the absence of chlorine, which results in low viscosity and (once cured) good weather resistance, low dielectric constants and high Tg. However, aliphatic epoxy resins polymerize very slowly at room temperature, so higher temperatures and suitable accelerators are usually required. Because aliphatic epoxies have a lower electron density than aromatics, cycloaliphatic epoxies react less readily with nucleophiles than bisphenol A-based epoxy resins (which have aromatic ether groups). This means that conventional nucleophilic hardeners such as amines are hardly suitable for crosslinking. Cycloaliphatic epoxides are therefore usually homopolymerized thermally or UV-initiated in an electrophilic or cationic reaction. Due to the low dielectric constants and the absence of chlorine, cycloaliphatic epoxides are often used to encapsulate electronic systems, such as microchips or LEDs. They are also used for radiation-cured paints and varnishes. Due to their high price, however, their use has so far been limited to such applications.
Epoxidized vegetable oils are formed by epoxidation of unsaturated fatty acids by reaction with peracids. In this case, the peracids can also be formed in situ by reacting carboxylic acids with hydrogen peroxide. Compared with LERs (liquid epoxy resins) they have very low viscosities. If, however, they are used in larger proportions as reactive diluents, this often leads to reduced chemical and thermal resistance and to poorer mechanical properties of the cured epoxides. Large scale epoxidized vegetable oils such as epoxidized soy and lens oils are used to a large extent as secondary plasticizers and cost stabilizers for PVC.
Aliphatic glycidyl epoxy resins of low molar mass (mono-, bi- or polyfunctional) are formed by the reaction of epichlorohydrin with aliphatic alcohols or polyols (glycidyl ethers are formed) or with aliphatic carboxylic acids (glycidyl esters are formed). The reaction is carried out in the presence of a base such as sodium hydroxide, analogous to the formation of bisphenol A-diglycidyl ether. Also aliphatic glycidyl epoxy resins usually have a low viscosity compared to aromatic epoxy resins. They are therefore added to other epoxy resins as reactive diluents or as adhesion promoters. Epoxy resins made of (long-chain) polyols are also added to improve tensile strength and impact strength.
A related class is cycloaliphatic epoxy resin, which contains one or more cycloaliphatic rings in the molecule (e.g. 3,4-epoxycyclohexylmethyl-3,4-epoxycyclohexane carboxylate). This class also displays lower viscosity at room temperature, but offers significantly higher temperature resistance than the aliphatic epoxy diluents. However, reactivity is rather low compared to other classes of epoxy resin, and high temperature curing using suitable accelerators is normally required. As aromaticity is not present in these materials as it is in Bisphenol A and F resins, the UV stability is considerably improved.
Halogenated
Halogenated epoxy resins are admixed for special properties, in particular brominated and fluorinated epoxy resins are used.
Brominated bisphenol A is used when flame retardant properties are required, such as in some electrical applications (e.g. printed circuit boards). The tetrabrominated bisphenol A (TBBPA, 2,2-bis(3,5-dibromophenyl)propane) or its diglycidyl ether, 2,2-bis[3,5-dibromo-4-(2,3-epoxypropoxy)phenyl]propane, can be added to the epoxy formulation. The formulation may then be reacted in the same way as pure bisphenol A. Some (non-crosslinked) epoxy resins with very high molar mass are added to engineering thermoplastics, again to achieve flame retardant properties.
Fluorinated epoxy resins have been investigated for some high performance applications, such as the fluorinated diglycidether 5-heptafluoropropyl-1,3-bis[2-(2,3-epoxypropoxy)hexafluoro-2-propyl]benzene. As it has a low surface tension, it is added as a wetting agent (surfactant) for contact with glass fibres. Its reactivity to hardeners is comparable to that of bisphenol A. When cured, the epoxy resin leads to a thermosetting plastic with high chemical resistance and low water absorption. However, the commercial use of fluorinated epoxy resins is limited by their high cost and low Tg.
Diluents
Epoxy resins diluents are typically formed by glycidylation of aliphatic alcohols or polyols and also aromatic alcohols. The resulting materials may be monofunctional (e.g. dodecanol glycidyl ether), difunctional (1,4-Butanediol diglycidyl ether), or higher functionality (e.g. trimethylolpropane triglycidyl ether). These resins typically display low viscosity at room temperature (10–200 mPa.s) and are often referred to as reactive diluents. They are rarely used alone, but are rather employed to modify (reduce) the viscosity of other epoxy resins. This has led to the term modified epoxy resin to denote those containing viscosity-lowering reactive diluents. The use of the diluent does effect mechanical properties and microstructure of epoxy resins. Mechanical properties of epoxy resins are generally not improved by use of diluents. Biobased epoxy diluents are also available.
Glycidylamine
Glycidylamine epoxy resins are higher functionality epoxies which are formed when aromatic amines are reacted with epichlorohydrin. Important industrial grades are triglycidyl-p-aminophenol (functionality 3) and N,N,N′,N′-tetraglycidyl-bis-(4-aminophenyl)-methane (functionality 4). The resins are low to medium viscosity at room temperature, which makes them easier to process than EPN or ECN resins. This coupled with high reactivity, plus high temperature resistance and mechanical properties of the resulting cured network makes them important materials for aerospace composite applications.
Curing
There are several dozen chemicals that can be used to cure epoxy, including amines, imidazoles, anhydrides and photosensitive chemicals. The study of epoxy curing is usually carried out by using differential scanning calorimetry.
In general, uncured epoxy resins have only poor mechanical, chemical and heat resistance properties. However, good properties are obtained by reacting the linear epoxy resin with suitable curatives to form three-dimensional cross-linked thermoset structures. This process is commonly referred to as curing or gelation process. Curing of epoxy resins is an exothermic reaction and in some cases produces sufficient heat to cause thermal degradation if not controlled. Curing does induce residual stress in epoxy systems which have been studied. The induced stresses may be alleviated with flexibilisers.
Curing may be achieved by reacting an epoxy with itself (homopolymerisation) or by forming a copolymer with polyfunctional curatives or hardeners. This curing is what produces the qualities of the substance such as resistance, durability, versatility, and adhesion. In principle, any molecule containing a reactive hydrogen may react with the epoxide groups of the epoxy resin. Common classes of hardeners for epoxy resins include amines, acids, acid anhydrides, phenols, alcohols and thiols. Relative reactivity (lowest first) is approximately in the order: phenol < anhydride < aromatic amine < cycloaliphatic amine < aliphatic amine < thiol.
While some epoxy resin/ hardener combinations will cure at ambient temperature, many require heat, with temperatures up to being common, and up to for some specialist systems. Insufficient heat during cure will result in a network with incomplete polymerisation, and thus reduced mechanical, chemical and heat resistance. Cure temperature should typically attain the glass transition temperature (Tg) of the fully cured network in order to achieve maximum properties. Temperature is sometimes increased in a step-wise fashion to control the rate of curing and prevent excessive heat build-up from the exothermic reaction.
Hardeners which show only low or limited reactivity at ambient temperature, but which react with epoxy resins at elevated temperature are referred to as latent hardeners. When using latent hardeners, the epoxy resin and hardener may be mixed and stored for some time prior to use, which is advantageous for many industrial processes. Very latent hardeners enable one-component (1K) products to be produced, whereby the resin and hardener are supplied pre-mixed to the end user and only require heat to initiate curing. One-component products generally have shorter shelf-lives than standard 2-component systems, and products may require cooled storage and transport.
The epoxy curing reaction may be accelerated by addition of small quantities of accelerators. Tertiary amines, carboxylic acids and alcohols (especially phenols) are effective accelerators. Bisphenol A is a highly effective and widely used accelerator, but is now increasingly replaced due to health concerns with this substance. The most widely used accelerator is 2,4,6-Tris(dimethylaminomethyl)phenol.
Homopolymerisation
Epoxy resin may be reacted with itself in the presence of an anionic catalyst (a Lewis base such as tertiary amines or imidazoles) or a cationic catalyst (a Lewis acid such as a boron trifluoride complex) to form a cured network. This process is known as catalytic homopolymerisation. The resulting network contains only ether bridges, and exhibits high thermal and chemical resistance, but is brittle and often requires elevated temperature for the curing process, so finds only niche applications industrially.
Epoxy homopolymerisation is often used when there is a requirement for UV curing, since cationic UV catalysts may be employed (e.g. for UV coatings).
Amines
Polyfunctional primary amines form an important class of epoxy hardeners. Primary amines undergo an addition reaction with the epoxide group to form a hydroxyl group and a secondary amine. The secondary amine can further react with an epoxide to form a tertiary amine and an additional hydroxyl group. Kinetic studies have shown the reactivity of the primary amine to be approximately double that of the secondary amine. Use of a difunctional or polyfunctional amine forms a three-dimensional cross-linked network.
Aliphatic, cycloaliphatic and aromatic amines are all employed as epoxy hardeners. Amine type hardeners will alter both the processing properties (viscosity, reactivity) and the final properties (mechanical, temperature and heat resistance) of the cured copolymer network. Thus amine structure is normally selected according to the application. Overall reactivity potential for different hardeners can roughly be ordered; aliphatic amines > cycloaliphatic amines > aromatic amines, though aliphatic amines with steric hindrance near the amino groups may react as slowly as some of the aromatic amines. Slower reactivity allows longer working times for processors. Temperature resistance generally increases in the same order, since aromatic amines form much more rigid structures than aliphatic amines. Aromatic amines were widely used as epoxy resin hardeners, due to the excellent end properties when mixed with a parent resin. Over the past few decades concern about the possible adverse health effects of many aromatic amines has led to increased use of aliphatic or cycloaliphatic amine alternatives. Amines are also blended, adducted and reacted to alter properties and these amine resins are more often used to cure epoxy resins than a pure amine such as TETA. Increasingly, water-based polyamines are also used to help reduce the toxicity profile among other reasons.
Anhydrides
Epoxy resins may be thermally cured with anhydrides to create polymers with significant property retention at elevated temperatures for extended periods of time. Reaction and subsequent crosslinking occur only after opening of the anhydride ring, e.g. by secondary hydroxyl groups in the epoxy resin. Homopolymerization may also occur between epoxide and hydroxyl groups. The high latency of anhydride hardeners makes them suitable for processing systems which require addition of mineral fillers prior to curing, e.g. for high voltage electrical insulators. Cure speed may be improved by matching anhydrides with suitable accelerators. For dianhydrides, and to a lesser extent, monoanhydrides, non-stoichiometric, empirical determinations are often used to optimize dosing levels. In some cases, blends of dianhydrides and monoanhydrides can improve metering and mixing with liquid epoxy resins.
Phenols
Polyphenols, such as bisphenol A or novolacs can react with epoxy resins at elevated temperatures (), normally in the presence of a catalyst. The resulting material has ether linkages and displays higher chemical and oxidation resistance than typically obtained by curing with amines or anhydrides. Since many novolacs are solids, this class of hardeners is often employed for powder coatings.
Thiols
Also known as mercaptans, thiols contain a sulfur which reacts very readily with the epoxide group, even at ambient or sub-ambient temperatures. While the resulting network does not typically display high temperature or chemical resistance, the high reactivity of the thiol group makes it useful for applications where heated curing is not possible, or very fast cure is required e.g. for domestic DIY adhesives and chemical rock bolt anchors. Thiols have a characteristic odour, which can be detected in many two-component household adhesives.
Applications
The applications for epoxy-based materials are extensive and they are considered very versatile. The applications include coatings, adhesives and composite materials such as those using carbon fiber and fiberglass reinforcements (although polyester, vinyl ester, and other thermosetting resins are also used for glass-reinforced plastic). The chemistry of epoxies and the range of commercially available variations allows cure polymers to be produced with a very broad range of properties. They have been extensively used with concrete and cementitious systems. In general, epoxies are known for their excellent adhesion, chemical and heat resistance, good-to-excellent mechanical properties and very good electrical insulating properties. Many properties of epoxies can be modified (for example silver-filled epoxies with good electrical conductivity are available, although epoxies are typically electrically insulating). Variations offering high thermal insulation, or thermal conductivity combined with high electrical resistance for electronics applications, are available.
As with other classes of thermoset polymer materials, blending different grades of epoxy resin, as well as use of additives, plasticizers or fillers is common to achieve the desired processing or final properties, or to reduce cost. Use of blending, additives and fillers is often referred to as formulating.
All quantities of mix generate their own heat because the reaction is exothermic. Large quantities will generate more heat and thus greatly increase the rate of the reaction and so reduce working time (pot-life). So it is good practice to mix smaller amounts which can be used quickly to avoid waste and to be safer. There are various methods of toughening them, as they can be brittle. Rubber toughening is a key technology used for toughening.
Paints and coatings
Two part epoxy coatings were developed for heavy duty service on metal substrates and use less energy than heat-cured powder coatings. These systems provide a tough, protective coating with excellent hardness. One part epoxy coatings are formulated as an emulsion in water, and can be cleaned up without solvents.
Epoxy coatings are often used in industrial and automotive applications since they are more heat resistant than latex-based and alkyd-based paints. Epoxy paints tend to deteriorate, known as "chalking out", due to UV exposure. Epoxy coatings have also been used in drinking water applications. Epoxy coatings find much use to protect mild and other steels due to their excellent protective properties.
Change in color, known as yellowing, is a common phenomenon for epoxy materials and is often of concern in art and conservation applications. Epoxy resins yellow with time, even when not exposed to UV radiation. Significant advances in understanding yellowing of epoxies were achieved by Down first in 1984 (natural dark aging) and later in 1986 (high-intensity light aging). Down investigated various room-temperature-cure epoxy resin adhesives suitable for use in glass conservation, testing their tendency to yellow. A fundamental molecular understanding of epoxy yellowing was achieved, when Krauklis and Echtermeyer discovered the mechanistic origin of yellowing in a commonly used amine epoxy resin, published in 2018. They found that the molecular reason for epoxy yellowing was a thermo-oxidative evolution of carbonyl groups in the polymeric carbon–carbon backbone via a nucleophilic radical attack.
Polyester epoxies are used as powder coatings for washers, driers and other "white goods". Fusion Bonded Epoxy Powder Coatings (FBE) are extensively used for corrosion protection of steel pipes and fittings used in the oil and gas industry, potable water transmission pipelines (steel), and concrete reinforcing rebar. Epoxy coatings are also widely used as primers to improve the adhesion of automotive and marine paints especially on metal surfaces where corrosion (rusting) resistance is important. Metal cans and containers are often coated with epoxy to prevent rusting, especially for foods like tomatoes that are acidic. Epoxy resins are also used for decorative flooring applications such as terrazzo flooring, chip flooring, and colored aggregate flooring.
Epoxies have been modified in a variety of ways, including reacting with fatty acids derived from oils to yield epoxy esters, which were cured the same way as alkyds. Typical ones were L8 (80% linseed) and D4 (40% dehydrated castor oil). These were often reacted with styrene to make styrenated epoxy esters, used as primers. Curing with phenolics to make drum linings, curing esters with amine resins and pre-curing epoxies with amino resins to make resistant top coats. Organic chains maybe used to hydrophobically modify epoxy resins and change their properties. The effect of chain length of the modifiers has been studied.
Adhesives
Epoxy adhesives are a major part of the class of adhesives called "structural adhesives" or "engineering adhesives" (that includes polyurethane, acrylic, cyanoacrylate, and other chemistries.) These high-performance adhesives are used in the construction of aircraft, automobiles, bicycles, boats, golf clubs, skis, snowboards, and other applications where high strength bonds are required. Epoxy adhesives can be developed to suit almost any application. They can be used as adhesives for wood, metal, glass, stone, and some plastics. They can be made flexible or rigid, transparent or opaque/colored, fast setting or slow setting. Epoxy adhesives are better in heat and chemical resistance than other common adhesives. In general, epoxy adhesives cured with heat will be more heat- and chemical-resistant than those cured at room temperature. The strength of epoxy adhesives is degraded at temperatures above .
Some epoxies are cured by exposure to ultraviolet light. Such epoxies are commonly used in optics, fiber optics, and optoelectronics.
Industrial tooling and composites
Epoxy systems are used in industrial tooling applications to produce molds, master models, laminates, castings, fixtures, and other industrial production aids. This "plastic tooling" replaces metal, wood and other traditional materials, and generally improves the efficiency and either lowers the overall cost or shortens the lead-time for many industrial processes. Epoxies are also used in producing fiber-reinforced or composite parts. They are more expensive than polyester resins and vinyl ester resins, but usually produce stronger and more temperature-resistant thermoset polymer matrix composite parts. Machine bedding to overcome vibrations is a use in the form of epoxy granite.
Wind turbine technology composites
Epoxy resins are used as bonding matrix along with glass or carbon fiber fabrics to produce composites with very high strength to weight characteristics, allowing longer and more efficient rotor blades to be produced. In addition, for offshore and onshore wind energy installations, epoxy resins are used as protective coatings on steel towers, base struts and concrete foundations. Aliphatic polyurethane top coats are applied on top to ensure full UV protection, prolong operational lifetimes and lowering maintenance costs. Electric generators, connected via the drivetrain with the rotor blades, convert mechanical wind energy to usable electric energy, and rely on epoxies electrical insulation and high thermal resistance properties. The same applies to transformers, bushings, spacers, and composites cables connecting the windmills to the grid. In Europe, wind energy components account for the largest segment of epoxy applications, about 27% of the market.
Electrical systems and electronics
Epoxy resin formulations are important in the electronics industry, and are employed in motors, generators, transformers, switchgear, bushings, insulators, printed wiring boards (PWB), and semiconductor encapsulants. Epoxy resins are excellent electrical insulators and protect electrical components from short circuiting, dust and moisture. In the electronics industry epoxy resins are the primary resin used in overmolding integrated circuits, transistors and hybrid circuits, and making printed circuit boards. The largest volume type of circuit board—an "FR-4 board"—is a sandwich of layers of glass cloth bonded into a composite by an epoxy resin. Epoxy resins are used to bond copper foil to circuit board substrates, and are a component of the solder mask on many circuit boards.
Flexible epoxy resins are used for potting transformers and inductors. By using vacuum impregnation on uncured epoxy, winding-to-winding, winding-to-core, and winding-to-insulator air voids are eliminated. The cured epoxy is an electrical insulator and a much better conductor of heat than air. Transformer and inductor hot spots are greatly reduced, giving the component a stable and longer life than unpotted product.
Epoxy resins are applied using the technology of resin dispensing.
Petroleum & petrochemical
Epoxies can be used to plug selective layers in a reservoir which are producing excessive brine. The technique is named "water shut-off treatment".
Consumer and marine applications
Epoxies are sold in hardware stores, typically as a pack containing separate resin and hardener, which must be mixed immediately before use. They are also sold in boat shops as repair resins for marine applications. Epoxies typically are not used in the outer layer of a boat because they deteriorate by exposure to UV light. They are often used during boat repair and assembly, and then over-coated with conventional or two-part polyurethane paint or marine-varnishes that provide UV protection.
There are two main areas of marine use. Because of the better mechanical properties relative to the more common polyester resins, epoxies are used for commercial manufacture of components where a high strength/weight ratio is required. The second area is that their strength, gap filling properties and excellent adhesion to many materials including timber have created a boom in amateur building projects including aircraft and boats.
Normal gelcoat formulated for use with polyester resins and vinylester resins does not adhere to epoxy surfaces, though epoxy adheres very well if applied to polyester resin surfaces. "Flocoat" that is normally used to coat the interior of polyester fibreglass yachts is also compatible with epoxies.
Epoxy materials tend to harden somewhat more gradually, while polyester materials tend to harden quickly, particularly if a lot of catalyst is used. The chemical reactions in both cases are exothermic.
While it is common to associate polyester resins and epoxy resins, their properties are sufficiently different that they are properly treated as distinct materials. Polyester resins are typically low strength unless used with a reinforcing material like glass fibre, are relatively brittle unless reinforced, and have low adhesion. Epoxies, by contrast, are inherently strong, somewhat flexible and have excellent adhesion. However, polyester resins are much cheaper.
Epoxy resins typically require a precise mix of two components which form a third chemical to get the stated properties. Depending on the properties required, the ratio may be anything from 1:1 or over 10:1, but in usually they must be mixed exactly. The final product is then a precise thermoset plastic. Until they are mixed the two elements are relatively inert, although the 'hardeners' tend to be more chemically active and should be protected from the atmosphere and moisture. The rate of the reaction can be changed by using different hardeners, which may change the nature of the final product, or by controlling the temperature.
By contrast, polyester resins are usually made available in a 'promoted' form, such that the progress of previously-mixed resins from liquid to solid is already underway, albeit very slowly. The only variable available to the user is to change the rate of this process using a catalyst, often Methyl-Ethyl-Ketone-Peroxide (MEKP), which is very toxic. The presence of the catalyst in the final product actually detracts from the desirable properties, so that small amounts of catalyst are preferable, so long as the hardening proceeds at an acceptable pace. The rate of cure of polyesters can therefore be controlled by the amount and type of catalyst as well as by the temperature.
As adhesives, epoxies bond in three ways: a) Mechanically, because the bonding surfaces are roughened; b) by proximity, because the cured resins are physically so close to the bonding surfaces that they are hard to separate; c) ionically, because the epoxy resins form ionic bonds at an atomic level with the bonding surfaces. This last is substantially the strongest of the three. By contrast, polyester resins can only bond using the first two of these, which greatly reduces their utility as adhesives and in marine repair.
Construction applications
Epoxies have been researched and used for construction for a few decades. Although they increase cost of mortars and concrete when used as an additive, they enhance properties. Research is ongoing to investigate the use of epoxies and other recycled plastics in mortars to enhance properties and recycle waste. Densifying plastic materials such as PET and plastic bags and then using them to partially replace aggregate and depolymerizing PET to use as a polymeric binder in addition to epoxy to enhance concrete are actively being studied.
Aerospace applications
In the aerospace industry, epoxy is used as a structural matrix material which is then reinforced by fiber. Typical fiber reinforcements include glass, carbon, Kevlar, and boron.
Epoxies are also used as a structural glue. Materials like wood, and others that are 'low-tech' are glued with epoxy resin. Epoxies generally out-perform most other resin types in terms of mechanical properties and resistance to environmental degradation.
Biology
Water-soluble epoxies such as Durcupan are commonly used for embedding electron microscope samples in plastic so they may be sectioned (sliced thin) with a microtome and then imaged.
Art
Epoxy resin, mixed with pigment, may be used as a painting medium, by pouring layers on top of each other to form a complete picture. It is also used in jewelry, as a doming resin for decorations and labels, and in decoupage type applications for art, countertops, and tables. Its seamless and glossy finish, along with its ability to be molded into various shapes, makes epoxy resin a favored choice for creating minimalist and statement pieces in modern furniture design and into various other design styles, including industrial, rustic, and even eclectic. It has been used and studied for art and historic structure preservation.
Production
The global epoxy resin market was valued at approximately $8 billion in 2016. The epoxy resin market is dominated by the Asia-Pacific region, which contributes 55.2% of the total market share. China is the major producer and consumer globally, consuming almost 35% of the global resin production. The global market is made up of approximately 50–100 manufacturers of basic or commodity epoxy resins and hardeners. In Europe, about 323,000 tonnes of epoxy resin were manufactured in 2017 generating some €1,055 million in sales. Germany is the largest market for epoxy resins in Europe, followed by Italy, France, the UK, Spain, the Netherlands and Austria.
These commodity epoxy manufacturers mentioned above typically do not sell epoxy resins in a form usable to smaller end users, so there is another group of companies that purchases epoxy raw materials from the major producers and then compounds (blends, modifies, or otherwise customizes) epoxy systems from these raw materials. These companies are known as "formulators". The majority of the epoxy systems sold are produced by these formulators and they comprise over 60% of the dollar value of the epoxy market. There are hundreds of ways that these formulators can modify epoxies—by adding mineral fillers (talc, silica, alumina, etc.), by adding flexibilizers, viscosity reducers, colorants, thickeners, accelerators, adhesion promoters, etc. These modifications are made to reduce costs, to improve performance, and to improve processing convenience. As a result, a typical formulator sells dozens or even thousands of formulations—each tailored to the requirements of a particular application or market.
The raw materials for epoxy resin production are today largely petroleum derived, although some plant derived sources are now becoming commercially available (e.g. plant derived glycerol used to make epichlorohydrin).
Renewable, recycled, waterborne and biobased epoxy
As there is a general trend to renewable and "green" sources and greater use of biobased materials, research is ongoing in the epoxy arena too. Waterborne epoxy paints have been around since the 1970s and research is ongoing. There is also movement to use waste as well as recycled raw materials where possible. Waste granite filings are produced in the mining industry. Research is being done on innovative solutions such as using waste granite powders in epoxy resins and designing binders for coatings based on this. Other work is ongoing to produce epoxy and epoxy based coatings from recycled raw materials including PET bottles. Monomers for epoxy resins can be derived from several different biomass sources.
Health risks
Liquid epoxy resins in their uncured state are mostly classed as irritant to the eyes and skin, as well as toxic to aquatic organisms. Solid epoxy resins are generally safer than liquid epoxy resins, and many are classified non-hazardous materials. One particular risk associated with epoxy resins is sensitization. The risk has been shown to be more pronounced in epoxy resins containing low molecular weight epoxy diluents. Exposure to epoxy resins can, over time, induce an allergic reaction. Sensitization generally occurs due to repeated exposure (e.g. through poor working hygiene or lack of protective equipment) over a long period of time. Allergic reaction sometimes occurs at a time which is delayed several days from the exposure. Allergic reaction is often visible in the form of dermatitis, particularly in areas where the exposure has been highest (commonly hands and forearms). Epoxy use is a main source of occupational asthma among users of plastics. Safe disposal also needs considering but usually involves deliberate curing to produce solid rather than liquid waste.
| Physical sciences | Polymers | Chemistry |
211761 | https://en.wikipedia.org/wiki/Mongoose | Mongoose | A mongoose is a small terrestrial carnivorous mammal belonging to the family Herpestidae. This family has two subfamilies, the Herpestinae and the Mungotinae. The Herpestinae comprises 23 living species that are native to southern Europe, Africa and Asia, whereas the Mungotinae comprises 11 species native to Africa.
The Herpestidae originated about in the Early Miocene and genetically diverged into two main lineages between 19.1 and .
Etymology
The name is derived from names used in India for Herpestes species:
or in classical Hindi;
in Marathi;
in Telugu;
, and in Kannada.
The form of the English name (since 1698) was altered to its "-goose" ending by folk etymology. It was spelled "mungoose" in the 18th and 19th centuries.
The plural form is "mongooses".
Characteristics
Mongooses have long faces and bodies, small, rounded ears, short legs, and long, tapering tails. Most are brindled or grizzly; a few have strongly marked coats which bear a striking resemblance to mustelids. Their nonretractile claws are used primarily for digging. Mongooses, much like goats, have narrow, ovular pupils. Most species have a large anal scent gland, used for territorial marking and signaling reproductive status and a short and smooth penis with a baculum and an elongated urethral opening on its underside. The dental formula of mongooses is .
They range from in head-to-body length, excluding the tail. In weight, they range from to .
Mongooses are one of at least four known mammalian taxa with mutations in the nicotinic acetylcholine receptor that protect against snake venom. Their modified receptors prevent the snake venom α-neurotoxin from binding. These represent four separate, independent mutations. In the mongoose, this change is effected, uniquely, by glycosylation.
Taxonomy
Herpestina was a scientific name proposed by Charles Lucien Bonaparte in 1845 who considered the mongooses a subfamily of the Viverridae. In 1864, John Edward Gray classified the mongooses into three subfamilies: Galidiinae, Herpestinae and Mungotinae. This grouping was supported by Reginald Innes Pocock in 1919, who referred to the family as "Mungotidae".
Genetic research based on nuclear and mitochondrial DNA analyses revealed that the Galidiinae are more closely related to Madagascar carnivores, including the fossa and Malagasy civet. Galidiinae is considered a subfamily of Eupleridae.
Phylogenetic relationships
Phylogenetic research of 18 mongoose species revealed that the solitary and social mongooses form different clades.
The phylogenetic relationships of Herpestidae are shown in the following cladogram:
Extinct species
Atilax
†A. mesotes
Herpestes
†H. lemanensis
Leptoplesictis
†L. atavus Beaumont, 1973
†L. aurelianensis Schlosser, 1888
†L. filholi Gaillard, 1899
†L. mbitensis Schmidt-Kittler, 1987
†L. namibiensis Morales et al., 2008
†L. peignei, Grohé et al., 2020
†L. rangwai Schmidt-Kittler, 1987
†L. senutae Morales et al., 2008
Behaviour and ecology
Mongooses mostly feed on insects, crabs, earthworms, lizards, birds, and rodents. However, they also eat eggs and carrion.
Some species can learn simple tricks. They can be tamed and are kept as pets to control vermin.
Cultural significance
In ancient Mesopotamia, mongooses were sacred to the deity Ninkilim, who was conflated with Ningirama, a deity of magic who was invoked for protection against serpents. According to a Babylonian popular saying, when a mouse fled from a mongoose into a serpent's hole, it announced, "I bring you greetings from the snake-charmer!" A creature resembling a mongoose also appears in Old Babylonian glyptic art, but its significance is not known.
All mongoose species, except for Suricata suricatta, are classed as a "prohibited new organism" under New Zealand's Hazardous Substances and New Organisms Act 1996, preventing them from being imported into the country.
A well-known fictional mongoose is Rikki-Tikki-Tavi, who appears in a short story of the same title in The Jungle Book (1894) by Rudyard Kipling. In this tale set in India, a young pet mongoose saves his human family from a krait and from Nag and Nagaina, two cobras. The story was later made into several films and a song by Donovan, among other references. A mongoose is also featured in Bram Stoker's novel The Lair of the White Worm. The main character, Adam Salton, purchases one to independently hunt snakes. Another mongoose features in the denouement of the Sherlock Holmes story "The Adventure of the Crooked Man", by Sir Arthur Conan Doyle. The Indian Tamil devotional film Padai Veetu Amman shows Tamil actor Vinu Chakravarthy changing himself into a mongoose by using his evil tantric mantra, to fight the goddess Amman. However, the mongoose finally dies at the hands of the goddess.
Mongoose species are prohibited to be kept as pets in the United States.
| Biology and health sciences | Carnivora | null |
211816 | https://en.wikipedia.org/wiki/Circinus | Circinus | Circinus is a small, faint constellation in the southern sky, first defined in 1756 by the French astronomer Nicolas-Louis de Lacaille. Its name is Latin for compass, referring to the drafting tool used for drawing circles (it should not be confused with Pyxis, a constellation that represents a mariner's compass which points north). Its brightest star is Alpha Circini, with an apparent magnitude of 3.19. Slightly variable, it is the brightest rapidly oscillating Ap star in the night sky. AX Circini is a Cepheid variable visible with the unaided eye, and BX Circini is a faint star thought to have been formed from the merger of two white dwarfs. Two sun-like stars have planetary systems: HD 134060 has two small planets, and HD 129445 has a Jupiter-like planet. Supernova SN 185 appeared in Circinus in 185 AD and was recorded by Chinese observers. Two novae have been observed more recently, in the 20th century.
The Milky Way runs through the constellation, featuring prominent objects such as the open cluster NGC 5823 and the planetary nebula NGC 5315. Circinus hosts a notable spiral galaxy, the Circinus Galaxy, discovered in 1977; it is the closest Seyfert galaxy to the Milky Way. The Alpha Circinids (ACI), a meteor shower also discovered in 1977, radiate from this constellation.
History
In 1756, French astronomer Nicolas-Louis de Lacaille introduced the constellation of Circinus with the French name le Compas, representing a pair of dividing compasses, on a chart of the southern sky. On that chart, Lacaille portrayed the constellations of Norma, Circinus, and Triangulum Australe, respectively, as a set square and ruler, a compass, and a surveyor's level in a set of draughtsman's instruments. Circinus was given its current name in 1763, when Lacaille published an updated sky map with Latin names for the constellations he introduced.
Characteristics
Bordered by Centaurus, Musca, Apus, Triangulum Australe, Norma and Lupus, Circinus lies adjacent to the Alpha and Beta Centauri stars. As it is at declination −50° to −70°, the whole constellation is only visible south of latitude 30° N. The official constellation boundaries, as set by Belgian astronomer Eugène Delporte in 1930, are defined by a polygon of 14 segments. In the equatorial coordinate system, the right ascension coordinates of these borders lie between and , and the declination coordinates are between −55.43° and −70.62°. Circinus culminates each year at 9 p.m. on 30 July. The recommended three-letter abbreviation for the constellation, as adopted by the International Astronomical Union in 1922, is "Cir".
Features
Stars
Circinus is a faint constellation, with only one star brighter than fourth magnitude. Alpha Circini, a white main sequence star with an apparent magnitude of 3.19, is 54 light-years away and 4° south of Alpha Centauri. Not only the brightest star in the constellation, it is also the brightest example of a rapidly oscillating Ap (RoAp) star in the night sky. It has the unusual spectral type A7 Vp SrCrEu, showing increased emissions of strontium, chromium and europium. Stars of this type have oddly localised magnetic fields and are slightly variable. Alpha Circini forms a binary star system with an orange dwarf companion of spectral type K5 and magnitude 8.5, which with a separation of 5.7 arcseconds is only discernible with a telescope. The distance between the two stars is 260 AU and they take 2600 years to rotate around a common centre of gravity. The second brightest star is Beta Circini, a white main sequence star of spectral type A3Va and a magnitude of 4.07, about 100 light-years away. It has around 1.8 times the diameter of the Sun.
Gamma Circini is a binary star 450 light-years away, whose components need a telescope of 150 mm to be seen, as they are only 0.8 arcseconds apart. The brighter component is a bluish Be star of spectral type B5IV+ and magnitude 4.51, while the dimmer component is a yellow star of magnitude 5.5. They orbit each other every 180 years. Delta Circini is also a multiple star whose components have magnitudes of 5.1 and 13.4 and orbit around a common centre of gravity every 3.9 days. The brighter component is a close eclipsing binary (specifically, a rotating ellipsoidal variable), with a minor dip of magnitude (0.1). Both are hot blue stars of spectral types O7III-V and O9.5V, respectively, and are estimated to have around 22 and 12 times the Sun's mass. Over 3600 light-years away, this system would outshine Venus at magnitude −4.8 if it were 32 light-years (10 parsecs) distant. The two main components are separated by 50 arcseconds, resolvable to the naked eye for individuals with good vision and easily discernible with a telescope.
Eta Circini is a yellow giant of spectral type G8III and magnitude 5.17, located around 276 light-years distant, and Zeta Circini is a blue-white main sequence star of spectral type B3V and magnitude 6.09, located around 1273 light-years away.
493 variable stars have been recorded in Circinus, but most have a very small range or are quite dim. Three prominent examples are Theta Circini, T Circini, and AX Circini. Theta Circini is a B-class irregular variable, ranging in magnitude from 5.0 to 5.4. T Circini has a B-type spectrum, ranging in magnitude from 10.6 to 9.3 over a period of 3.298 days, although it is actually an eclipsing binary system rather than a pulsating star. AX is a Cepheid variable that varies between magnitudes 5.6 and 6.19 over 5.3 days. It is a yellow-white supergiant of spectral type F8II+, 1600 light-years away. BP Circini is another Cepheid variable with an apparent magnitude ranging from 7.37 to 7.71 over 2.4 days. Both cepheids are spectroscopic binaries, with companions that are blue-white stars of spectral type B6 and 5 and 4.7 solar masses, respectively. BX Circini is a faint star that fluctuates between magnitudes 12.57 and 12.62 over a period of 2 hours 33 minutes. Over 99% of its composition appears to be helium. Its origin is unclear, but thought to be the result of the merger of a helium and a carbon/oxygen white dwarf.
Several stars with planetary systems lie within the borders of Circinus, although none of the host stars are particularly prominent. HD 134060 is a sun-like yellow dwarf star of spectral type G0VFe+0.4 and magnitude 6.29, around 79 light-years away. Its two planets were discovered in 2011 through the radial velocity method: the smaller, HD 134060 b, has a mass of 0.0351 MJ (Jupiter masses) and orbits its star every 3.27 days, at 0.0444 AU; the larger, HD 134060 c (0.15 MJ), orbits farther out at 2.226 AU, with a period of approximately 1161 days. Even fainter, at magnitude 8.8, HD 129445 is 220 light-years away and has 99% of the Sun's mass and a similar spectral type of G8V. HD 129445 b, a Jupiter-like planet (1.6 MJ) discovered in 2010 via the radial velocity method, orbits this star at a distance of 2.9 AU, approximately every 1840 days.
As this constellation intersects the plane of the Milky Way, there are many massive stars located in this constellation, including GKF2010 MN18 (or simply MN18), a blue supergiant located in a bipolar nebula, as well as 9 Wolf-Rayet stars, a very high number for such a small constellation.
Deep-sky objects
Three open clusters and a planetary nebulae are found within the borders of Circinus, all visible with amateur telescopes of varying sizes. NGC 5823, also called Caldwell 88, is an 800-million-year-old open cluster, located 3500 light-years away and spanning a 12-light-year region along the constellation's northern border. Despite having an integrated magnitude of 7.9, the cluster can be seen by star hopping from Beta Circini or from Alpha Centauri. It contains 80–100 stars of 10th magnitude and fainter, which are spread out over a diameter of 10 arcseconds. The brighter stars, however, are not true members of the cluster, as they are closer to the Earth than the dimmer ones. NGC 5823 appears distinct to the observer, sometimes seen as a reversed "S", as described by John Herschel, although it has also been described as "tulip-shaped" and "boxy". That cluster can be easily mistaken with a similar cluster, NGC 5822, nearby in Lupus. Comparatively, open cluster NGC 5715 is fainter (integrated magnitude of 9.8)—its brightest star is only 11th magnitude—and smaller (7.0 arcminutes), comprising only 30 stars. The third open cluster, Pismis 20, contains 12 stars in a diameter of 4.5 arcseconds but exhibits a magnitude similar to NGC 5823 (7.8). At 8270 light-years, it requires an amateur telescope with an aperture over 300 mm to be easily discerned.
The planetary nebula NGC 5315 has a magnitude of 9.8 around a central star of magnitude 14.2, located 5.2 degrees west-southwest of Alpha Circini. It is only visible as a disc at magnifications over 200-fold. Bernes 145 is a dark and reflection nebula first listed in the 1971 Bernes Catalog. The dark nebula component is easily visible in a large amateur telescope, and it measures 12 by 5 arcminutes. The smaller reflection nebula component requires a larger instrument and averted vision to be seen.
Circinus also houses ESO 97-G13, commonly known as the Circinus Galaxy. Discovered in 1977, it is a relatively unobscured galaxy (magnitude 10.6), which is unusual for galaxies located in constellations near the Milky Way, since their dim light is obscured by gas and dust. This oblong spiral galaxy with 6.9 by 3.0 arcminutes and 26,000 light-years in diameter, is located 13 million light-years away from Earth and lies 4 degrees off the galactic plane. It is the closest Seyfert galaxy to the Milky Way, and therefore hosts an active galactic nucleus.
Circinus X-1 is an X-ray binary star system that includes a neutron star. Observations of Circinus X-1 in July 2007 revealed the presence of X-ray jets normally found in black hole systems. Located at 19,000 light-years, the pulsar PSR B1509-58, also called the Circinus Pulsar, has expelled a 20-light-year-long jet of material from its southern pole, clearly visible in the X-ray spectrum. Another supernova remnant in Circinus is that of SN 185. Recorded by Chinese observers in 185 AD, SN 185 was visible in the night sky for around eight months; its remnants, known as RCW 86, cover an area larger than the typical full moon.
A white dwarf star in a close binary system can accumulate material from its companion until it ignites and blows off in a thermonuclear explosion, known as a nova. These stars generally brighten by 7 to 16 magnitudes. Nova Circini 1926, also known as X Circini, was observed at magnitude 6.5 on 3 September 1926, before fading and fluctuating between magnitudes 11.7 and 12.5, during 1928, and magnitude 13, in 1929. Nova Circini 1995 (BY Circini) reached a maximum apparent magnitude of 7.2 in January 1995. BW Circini is a low mass X-ray binary system, comprising a black hole of around 8 solar masses and a yellow G0III-G5III subgiant star. X-ray outbursts were recorded in 1987 and 1997, and possibly 1971–72.
Meteor showers
Circinus is the radiant of an annual meteor shower, the Alpha Circinids (ACI). First observed in Queensland in 1977, the meteors have an average velocity of 27.1 km/s and are thought to be associated with a long-period comet. In 2011, Peter Jenniskens proposed that the debris trail of comet C/1969 T1 could intersect with the Earth's orbit and generate a meteor outburst coming from a radiant close to Beta Circini. The ACI shower peaks on 4 June, the day it was first observed.
| Physical sciences | Other | Astronomy |
211818 | https://en.wikipedia.org/wiki/Volans | Volans | Volans is a constellation in the southern sky. It represents a flying fish; its name is a shortened form of its original name, Piscis Volans. Volans was one of twelve constellations created by Petrus Plancius from the observations of Pieter Dirkszoon Keyser and Frederick de Houtman and it first appeared on a 35-cm (14") diameter celestial globe published in 1597 (or 1598) in Amsterdam by Plancius with Jodocus Hondius. The first depiction of this constellation in a celestial atlas was in Johann Bayer's Uranometria of 1603.
History
Volans is one of the 12 constellations that were introduced by the Dutch navigators Pieter Dirkszoon Keyser and Frederick de Houtman in the late 16th century. It was first depicted on Petrus Plancius’ globe in 1598. Plancius called the constellation Vliegendenvis (flying fish).
In 1603, Johann Bayer included the constellation in his star atlas Uranometria under the name Piscis Volans, the flying fish. John Herschel proposed shrinking the name to one word in 1844, noting that Lacaille himself had abbreviated his constellations thus on occasion. This was universally adopted.
Volans represents a type of tropical fish that can jump out of the water and glide through the air on wings. In early celestial maps, the flying fish was often depicted as accompanying the ship Argo Navis, and being chased by the predatory fish represented by the adjoining constellation Dorado.
On 10 May 2023, TOI-715 b, the first exoplanet in the conservative habitable zone, about 1.55 times larger than Earth, was discovered in Volans by TESS.
Features
Stars
There are two double stars within the constellation which can be observed using a small telescope, Gamma Volantis and Epsilon Volantis, along with two galaxies which may be more difficult to see clearly, NGC 2442 and NGC 2434. The magnitudes of the Gamma Volantis stars are 3.8 and 5.6, and of Epsilon Volantis 4.4 and 7.3.
The Astronomical Society of Southern Africa in 2003 reported that observations of the variable stars R and S Volantis in Volans were very urgently needed as data on their light curves was incomplete.
HD 76700 is a sunlike star some 195 light-years distant that has been found to have a planet.
Deep-sky objects
Volans has several deep-sky objects within its borders.
The Lindsay-Shapley ring, also categorized as AM0644-741, is a ring galaxy located 300 million light-years from Earth. Named for its discoverers, the Lindsay-Shapley ring was found near the Large Magellanic Cloud in 1960. Like the Cartwheel Galaxy in Sculptor, the unusual shape of this galaxy results from a collision many millions of years ago. The blue ring, 150,000 light-years in diameter, was formed when a shock wave from the collision created a ring of hot blue stars; the yellow core is an amalgamation of the progenitors' cores. NGC 2442, an intermediate-spiral galaxy, is also located in this constellation, with a distance of 50 million light-years from Earth.
Graham's Object, aka Das Rheingold or Nibelungen Ring is a ring shaped galaxy at R.A. 6h 41.4m / Decl. -74° 19' (2000.0) in Volans
| Physical sciences | Other | Astronomy |
211819 | https://en.wikipedia.org/wiki/Reticulum | Reticulum | Reticulum is a small, faint constellation in the southern sky. Its name is Latin for a small net, or reticle—a net of crosshairs at the focus of a telescope eyepiece that is used to measure star positions. The constellation is best viewed between October and December, and save for one main star visible in ideal conditions, cannot be seen from north of the 30th parallel north.
History
A constellation in this area was introduced by Isaac Habrecht II in his celestial globe in 1621, who named it Rhombus. It was replaced with a somewhat different constellation by the French astronomer Nicolas Louis de Lacaille in the eighteenth century; during his stay at the Cape of Good Hope, he named the constellation le Réticule Rhomboide to commemorate the reticle in his telescope eyepiece. The name was later Latinized to Reticulum in his star catalogue Coelum Australe Stelliferum. In 1810, the stars of Reticulum were used by William Croswell to produce the constellation Marmor Sculptile, which represented the bust of Christopher Columbus, but this did not catch on among astronomers.
The constellation Reticulum became officially recognized during the First General Assembly of the International Astronomical Union in 1922. The boundary for this and other constellations was drawn up by Belgian astronomer Eugène Delporte along arcs of right ascension and declination for epoch 1875. These were published in 1930 in the Delimination Scientifique des Constellations at the behest of the IAU.
Features
Only two of the stars in this constellation are brighter than visual magnitude 5: Alpha (α) and Beta (β) reticuli. The reddish star R Reticuli is a Mira variable. This variable was discovered by C. Ragoonatha Chary at the Madras Observatory in India.
The binary star system Epsilon Reticuli consists of a spectral class K2IV star being orbited by a white dwarf. Based on parallax measurements, this system is located about 50 light years from the Sun. In 2000, a planetary companion was announced, orbiting the star ε Reticuli A.
Zeta Reticuli is a wide binary star system, with both members being similar to the Sun. It is located at a distance of about 39 light years. This system gained some notoriety in ufology when the alleged alien abductees Betty and Barney Hill named it as the home of their abductors.
In 2005, a type Ia supernova was discovered in the spiral galaxy NGC 1559, located in the constellation.
The dwarf galaxy Reticulum II is enriched in r-process heavy elements.
The Horologium-Reticulum Supercluster is a galaxy supercluster that ranges from 700 million to 1.2 billion light-years from Earth.
| Physical sciences | Other | Astronomy |
211820 | https://en.wikipedia.org/wiki/Pictor | Pictor | Pictor is a constellation in the Southern Celestial Hemisphere, located between the star Canopus and the Large Magellanic Cloud. Its name is Latin for painter, and is an abbreviation of the older name Equuleus Pictoris (the "painter's easel"). Normally represented as an easel, Pictor was named by Abbé Nicolas-Louis de Lacaille in the 18th century. The constellation's brightest star is Alpha Pictoris, a white main-sequence star around 97 light-years away from Earth. Pictor also hosts RR Pictoris, a cataclysmic variable star system that flared up as a nova, reaching apparent (visual) magnitude 1.2 in 1925 before fading into obscurity.
Pictor has attracted attention because of its second-brightest star Beta Pictoris, 63.4 light-years distant from Earth, which is surrounded by an unusual dust disk rich in carbon, as well as two exoplanets (extrasolar planets). Another five stars in the constellation have been observed to have planets. Among them is HD 40307, an orange dwarf that has six planets orbiting it, one of which—HD 40307 g—is a potential super-Earth in the circumstellar habitable zone. Kapteyn's Star, the nearest star in Pictor to Earth, is a red dwarf located 12.76 light-years away that was believed to have two super-Earths in orbit in 2014, but their existence of these planets was disproven in 2021. Pictor A is a radio galaxy that is shooting an 800,000 light-year long jet of plasma from a supermassive black hole at its centre. In 2006, a gamma-ray burst—GRB 060729—was observed in Pictor, its extremely long X-ray afterglow detectable for nearly two years.
History
The French astronomer Abbé Nicolas-Louis de Lacaille first described Pictor as le Chevalet et la Palette (the easel and palette) in 1756, after observing and cataloguing 10,000 southern stars during a two-year stay at the Cape of Good Hope. He devised 14 new constellations in uncharted regions of the Southern Celestial Hemisphere not visible from Europe. All but one honored instruments that symbolised the Age of Enlightenment. He gave these constellations Bayer designations, including ten stars in Pictor now named Alpha to Nu Pictoris. He labelled the constellation Equuleus Pictorius on his 1763 chart, the word "Equuleus" meaning small horse, or easel—perhaps from an old custom among artists of carrying a canvas on a donkey. The German astronomer Johann Bode called it Pluteum Pictoris. The name was shortened to its current form in 1845 by the English astronomer Francis Baily on the suggestion of his countryman Sir John Herschel.
Characteristics
Pictor is a small constellation bordered by Columba to the north, Puppis and Carina to the east, Caelum to the northwest, Dorado to the southwest and Volans to the south. The three-letter abbreviation for the constellation, as adopted by the International Astronomical Union in 1922, is "Pic". The official constellation boundaries, as set by Belgian astronomer Eugène Delporte in 1930, are defined by a polygon of 18 segments (illustrated in infobox). In the equatorial coordinate system, the right ascension coordinates of these borders lie between and , while the declination coordinates are between −42.79° and −64.15°. Pictor culminates each year at 9 p.m. on 17 March. Its position in the far Southern Celestial Hemisphere means that the whole constellation is visible to observers south of latitude 26°N, and parts become circumpolar south of latitude 35°S.
Features
Stars
Pictor is a faint constellation; its three brightest stars can be seen near the prominent Canopus. Within the constellation's borders, there are 49 stars brighter than or equal to apparent magnitude 6.5. Located about 97 light-years away from Earth, Alpha Pictoris is the brightest star in the constellation; it is a white main-sequence star with an apparent magnitude of 3.3, and spectral type A8VnkA6. A rapidly spinning star with a projected rotational velocity estimated at 206 km/s, it has a shell of circumstellar gas. Beta Pictoris is another white main sequence star of spectral type A6V and apparent magnitude 3.86. Located around 63.4 light-years distant from Earth, it is a member of the Beta Pictoris moving group—a group of 17 star systems around 12 million years old moving through space together. In 1984 Beta Pictoris was the first star discovered to have a debris disk. Since then, an exoplanet about eight times the mass of Jupiter has been discovered orbiting approximately 8 astronomical units (AU) away from the star—a similar distance as that between the Sun and Saturn. The European Southern Observatory (ESO) confirmed its presence through the use of direct imagery with the Very Large Telescope in late 2009.
Gamma Pictoris is an orange giant of spectral type K1III that has swollen to 1.4 times the diameter of the Sun. Shining with an apparent magnitude of 4.5, it lies 174 light-years distant from Earth. HD 42540, called 47 Pictoris by American astronomer Benjamin Apthorp Gould, is a slightly cooler orange giant, with a spectral type of K2.5III and average magnitude 5.04. It has also been suspected of being a variable star. Lacaille mistakenly named this star Mu Doradus, but had recorded its Right Ascension one hour too low. Lacaille named two neighbouring stars Eta Pictoris. Eta2 Pictoris, also known as HR 1663, is an orange giant of spectral type K5III and apparent magnitude 5.05. 474 light-years distant, it has a diameter 5.6 times that of the Sun. Eta1 Pictoris, also known as HR 1649, is 85 light-years distant and is a main sequence star of spectral type F5V and visual magnitude 5.38. A double star, it has a companion of magnitude 13; the two are separated by 11 arcseconds.
Located about 1298 light-years from Earth, Delta Pictoris is an eclipsing binary of the Beta Lyrae type. Composed of two blue stars of spectral types B3III and O9V, the system has a period of 1.67 days, and is observed to dip from apparent magnitude 4.65 to 4.9. The stars are oval-shaped as they are gravitationally distorted by each other. TV Pictoris is a spectroscopic binary system composed of an A-type star and an F-type star which rotate around each other in a very close orbit. The latter star is elliptical in shape and itself varies in brightness. The visual magnitude ranges between 7.37 and 7.53 every 20 hours.
Aside from Beta, five other stars in Pictor are known to host planetary systems. AB Pictoris is a BY Draconis variable star with a substellar companion that is either a large planet or a brown dwarf, which was discovered by direct imaging in 2005. HD 40307 is an orange main sequence star of spectral type K2.5V and apparent magnitude 7.17 located about 42 light-years away. Doppler spectroscopy with the High Accuracy Radial Velocity Planet Searcher (HARPS) indicates that HD 40307 is host to six super-Earth planets, one of which, HD 40307 g, lies in the circumstellar habitable zone of the star, and is not close enough to be tidally locked (i.e. with the same face always facing the star), unlike the other planets in the same system, and many other planets which orbit close to their parent stars. HD 41004 is a complex binary system about 139 light-years distant. The primary is an orange dwarf of spectral type K1V orbited by a planet roughly 2.65 times the mass of Jupiter every 963 days, while the secondary is a red dwarf of spectral type M2V and orbited by a brown dwarf that is at least 19 times as massive as Jupiter. Both substellar components were discovered by doppler spectroscopy using the CORALIE spectrograph in 2004 and 2002 respectively. Kapteyn's Star, a nearby red dwarf at the distance of 12.78 light-years, has a magnitude of 8.8. It has the largest proper motion of any star in the sky after Barnard's Star. Moving around the Milky Way in the opposite direction to most other stars, it may have originated in a dwarf galaxy that was merged into the Milky Way, with the main remnant being the Omega Centauri globular cluster. In 2014 analysis of the doppler variations of Kapteyn's Star with the HARPS spectrograph showed that it hosts two super-Earths—Kapteyn b and Kapteyn c, but the existence of these exoplanet was disproven in 2021. It is believed that these planets were actually just artifacts of the Kapteyn' star's rotation and activity.
Located 1.5 degrees west southwest of Alpha, RR Pictoris is a cataclysmic variable that flared up as a nova, reaching magnitude 1.2 on 9 June 1925. Six months after its peak brightness, it had faded to be invisible to the unaided eye, and was magnitude 12.5 by 1975. RR Pictoris is a close binary system composed of a white dwarf and secondary star that orbit each other every 3.48 hours—so close that the secondary is filling up its Roche lobe with stellar material, which is then transferred onto the first star's accretion disk. Once this material reaches a critical mass, it ignites and the system brightens tremendously. Calculations from the orbital speed suggest the secondary star is not dense enough for its size to still be on the main sequence, so it also must have begun expanding and cooling already after its core ran out of hydrogen fuel. The RR Pictoris system is estimated to lie around 1300 light-years distant from Earth.
Deep-sky objects
NGC 1705 is an irregular dwarf galaxy 17 million light-years from Earth. It is one of the most active star forming galaxies in the nearby universe, despite the fact that its rate of star formation peaked around 30 million years ago. Pictor A, around 485 million light-years away, is a double-lobed radio galaxy and a powerful source of radio waves in the Southern Celestial Hemisphere. From a supermassive black hole at its centre, a relativistic jet shoots out to an X-ray hot spot 800,000 light years away. SPT-CL J0546-5345 is a massive galaxy cluster located around 7 billion light-years away with a mass equivalent to approximately 800 trillion suns.
GRB 060729 was a gamma-ray burst that was first observed on 29 July 2006. It is likely the signal of a type Ic supernova—the core collapse of a massive star. It was also notable for its extraordinarily long X-ray afterglow, detectable 642 days (nearly two years) after the original event. The event was remote, with a redshift of 0.54.
| Physical sciences | Other | Astronomy |
211821 | https://en.wikipedia.org/wiki/Lacerta | Lacerta | Lacerta is one of the 88 modern constellations defined by the International Astronomical Union. Its name is Latin for lizard. A small, faint constellation, it was defined in 1687 by the astronomer Johannes Hevelius. Its brightest stars form a "W" shape similar to that of Cassiopeia, and it is thus sometimes referred to as 'Little Cassiopeia'. It is located between Cygnus, Cassiopeia and Andromeda on the northern celestial sphere. The northern part lies on the Milky Way.
Notable features
Lacerta is typical of Milky Way constellations: no bright galaxies, nor globular clusters, but instead open clusters, for example NGC 7243, the faint planetary nebula IC 5217 and quite a few double stars. It also contains the prototypic blazar BL Lacertae.
Lacerta contains no Messier objects.
Stars
Alpha Lacertae is a blue-white hued main-sequence star of magnitude 3.8, 102 light-years from Earth. It has a spectral type of A1 V and is an optical double star. Beta Lacertae is far dimmer, a yellow giant of magnitude 4.4, 170 light-years from Earth.
Roe 47 is a multiple star consisting of five components (magnitudes 5.8, 9.8, 10.1, 9.4, 9.8).
ADS 16402 is a binary star system in Lacerta, around which a planet orbits with some unusual properties. The Jupiter-sized planet exhibits an unexpectedly low density, about the same as cork. This planet is dubbed HAT P-1.
EV Lacertae is a rapidly spinning magnitude 10 red dwarf with a strong magnetic field. It is a flare star that can emit powerful flares potentially visible to the naked eye, thousands of times more energetic than any from Earth's sun.
HD 215441 (22h 44m.2 / +55° 35') is known as Babcock's Magnetic Star. Source: Sky Catalogue 2000.0 Volume 2: Double Stars, Variable Stars and Nonstellar Objects, page xlv, chapter Glossary of Selected Astronomical Names.
Deep-sky objects
NGC 7243 is an open cluster 2500 light-years from Earth, visible in small amateur telescopes. It has a few dozen "scattered" stars, the brightest of which are of the 8th magnitude.
BL Lacertae is the prototype of the BL Lacertae objects, which appear to be dim variable stars but are actually the variable nuclei of elliptical galaxies; they are similar to quasars. It lent its name to a whole type of celestial objects, the BL Lacertae objects (a subtype of blazar). The object varies irregularly between magnitudes 14 and 17 over a few days.
History
Centred on a region of the sky without apparently bright stars, Lacerta was apparently not regarded as a constellation by ancient Western astronomers. Johannes Hevelius created the constellation in 1687 and initially named it "Stellio" (the stellion), a lizard with star-like dorsal spots found along the Mediterranean coast.
Other Europeans who sought to name this new constellation included Augustin Royer, who created Sceptrum et Manus Iustitiae (the Hand of Justice and Sceptre) to honor Louis XIV in 1670, and Johann Elert Bode, who created Frederici Honores (Frederick's Glory) to honor Frederick the Great in 1787. Both Sceptrum and Frederici Honores are now obsolete, while Lacerta still survives.
Equivalents
Although not included in ancient star charts of Europe and the Near East, the stars of Lacerta, along with some in the eastern portion of Cygnus, were coincidentally combined by early Chinese astronomers into their "Flying serpent". Similarly, the Chumash people of California call this part of the sky 'Lizard' and include it in multiple stories.
Namesakes
USS Lacerta (AKA-29) was an attack cargo ship in the United States navy named after the constellation.
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211825 | https://en.wikipedia.org/wiki/Musca | Musca | is a small constellation in the deep southern sky. It was one of 12 constellations created by Petrus Plancius from the observations of Pieter Dirkszoon Keyser and Frederick de Houtman, and it first appeared on a celestial globe in diameter published in 1597 (or 1598) in Amsterdam by Plancius and Jodocus Hondius. The first depiction of this constellation in a celestial atlas was in Johann Bayer's Uranometria of 1603. It was also known as for 200 years. Musca remains below the horizon for most Northern Hemisphere observers.
Many of the constellation's brighter stars are members of the Scorpius–Centaurus association, a loose group of hot blue-white stars that appears to share a common origin and motion across the Milky Way. These include Alpha, Beta, Gamma, Zeta2 and (probably) Eta Muscae, as well as HD 100546, a blue-white Herbig Ae/Be star that is surrounded by a complex debris disk containing a large planet or brown dwarf and possible protoplanet. Two further star systems have been found to have planets. The constellation also contains two cepheid variables visible to the naked eye. Theta Muscae is a triple star system, the brightest member of which is a Wolf–Rayet star.
History
Musca was one of the 12 constellations established by the astronomer Petrus Plancius from the observations of the southern sky by the Dutch explorers Pieter Dirkszoon Keyser and Frederick de Houtman, who had sailed on the first Dutch trading expedition, known as the Eerste Schipvaart, to the East Indies. De Houtman included it in his southern star catalogue in 1598 under the Dutch name De Vlieghe, "The Fly". They assigned four stars to the constellation, with a star that would be later designated as Beta Muscae marking the head, Gamma the body, and Alpha and Delta the left and right wings, respectively. It first appeared on a 35-cm-diameter (14-in) celestial globe published in 1598 in Amsterdam by Plancius with Jodocus Hondius, though was unnamed. The first depiction of this constellation in a celestial atlas was in the German cartographer Johann Bayer's Uranometria of 1603, though Bayer termed it Apis— "the Bee", a name by which it was known for the next two centuries. A 1603 celestial globe by Willem Blaeu depicts it as providing nourishment for the nearby constellation Chamaeleon—its tongue trying to catch the insect.
The French explorer and astronomer Nicolas Louis de Lacaille called it la Mouche on the 1756 version of his planisphere of the southern skies. Jean Fortin retained the French name in 1776 for his Atlas Céleste, while Lacaille latinised the name for his revised Coelum Australe Stelliferum in 1763. Lacaille renamed it to Musca Australis, the Southern Fly—Australis, since it counterparted the now discarded constellation of Musca Borealis composed of a few stars in Aries, and to avoid confusion with Apus. Today, the name is simply Musca. It is the only official constellation depicting an insect.
The Kalapalo people of Mato Grosso state in Brazil called Alpha and Beta Muscae (along with Beta and Kappa Crucis) Kutsu anangagï "Ornate Hawk-Eagle's double flutes". The Wardaman people of the Northern Territory in Australia perceived the main stars of Musca as a ceremonial boomerang, part of the Central Arena—a sacred area surrounding the constellation Crux that depicts the lightning creation beings and where they teach Wardaman customs; Alpha and Beta also signified a ceremonial headband, while Gamma and Delta represented two armbands. In central Australia, the Arrernte and Luritja peoples living in on a mission in Hermannsburg viewed the sky as divided between them, east of the Milky Way representing Arrernte camps and west denoting Luritja camps. The stars of Musca, along with Fomalhaut, Alpha Pavonis, and Alpha and Beta Gruis, were all claimed by the Arrernte.
Characteristics
Musca is bordered by Crux to the north, Carina to the west, Chamaeleon to the south, Apus and Circinus to the east, and Centaurus to the northeast. Covering 138 square degrees and 0.335% of the night sky, it ranks 77th of the 88 constellations in size. The three-letter abbreviation for the constellation, as adopted by the International Astronomical Union in 1922, is "Mus". The official constellation boundaries, as set by Belgian astronomer Eugène Delporte in 1930, are defined by a polygon of six segments. In the equatorial coordinate system, the right ascension coordinates of these borders lie between and , while the declination coordinates are between −64.64° and −75.68°. The whole constellation is visible to observers south of latitude 14°N.
Features
Stars
Lacaille charted and designated 10 stars with the Bayer designations Alpha to Kappa in 1756. He catalogued stars that became Lambda and Mu, but did not designate them as he considered them informes as they lay outside the asterism proper. Baily considered them part of Musca, and Gould gave them their Bayer designations. Francis Baily also dropped Kappa, which he felt was too faint to warrant a name, and designated two adjacent stars as Zeta1 and Zeta2. These last two stars are 1° apart, quite far to be sharing a Bayer designation. Lacaille had originally labelled the fainter one as Zeta, while Baily presupposed he had meant to label the brighter one. Reluctant to remove Lacaille's designation, he gave them both the Zeta designation. Altogether there are 62 stars brighter than magnitude 6.5 in the constellation.
The pattern of the brightest stars resembles that of Ursa Minor, in that the stars form a pattern reminiscent of a bowl with a handle. Lying south-southeast of Acrux in neighbouring Crux is Alpha Muscae. It is the brightest star in the constellation with an apparent magnitude of 2.7. Lying around 310 light-years away, it is a blue-white star of spectral type B2IV-V that is around 4520 times as luminous and 8 times as massive as the Sun. The star is a Beta Cephei variable with about 4.7 times the Sun's diameter, and pulsates every 2.2 hours, varying by 1% in brightness. A nearby star of magnitude 13 may or may not be a companion star. Marking the fly's tail is Gamma Muscae, a blue-white star of spectral type B5V that varies between magnitudes 3.84 and 3.86 over a period of 2.7 days. It is a variable of a different type, classed as a slowly pulsating B star, a type of variable. It is around five times as massive as the Sun.
Beta Muscae is a binary star system around 341 light-years distant that is composed of two blue-white main-sequence stars of spectral types B2V and B3V that orbit each other every 194 years. They are eight and six times as massive as the Sun, respectively, and have about 3.5 times its diameter. Zeta2 Muscae is a white main sequence star of spectral type A5V around 330 light-years distant from Earth. It is part of a triple star system with faint companions at 0.5 and 32.4 arc seconds distance. Eta Muscae is a multiple star system, the two main components forming an eclipsing binary that has a combined spectral type of B8V and magnitude of 4.77 that dips by 0.05 magnitude every 2.39 days. Alpha, Beta, Gamma, HD 103079, Zeta2 and (likely) Eta are all members of the Lower Centaurus Crux subgroup of the Scorpius–Centaurus association, a group of predominantly hot blue-white stars that share a common origin and proper motion across the galaxy.
Delta and Epsilon mark the fly's left wing and right wing, respectively. With an apparent magnitude of 3.62, Delta is an orange giant of spectral type K2III located around 91 light-years away. Epsilon Muscae is a red giant of spectral type M5III and semiregular variable that ranges between magnitudes 3.99 and 4.31 over approximately 40 days. It has expanded to 130 times the Sun's diameter and 1800 to 2300 its luminosity. It was a star originally 1.5 to 2 times as massive as the Sun. Although of a similar distance—around 302 light-years—to the stars of the Lower Centaurus Crux subgroup, it is moving much faster at around 100 km/s and does not share a common origin. To the northwest lies Mu Muscae, an orange giant of spectral type K4III that varies between apparent magnitude 4.71 and 4.76, and has been classified as a slow, irregular variable. Near Mu is Lambda Muscae, the third-brightest star in the constellation and a white main-sequence star of spectral type A7V around 128 light-years distant from Earth.
Located near Alpha is R Muscae, a classical Cepheid variable ranging from apparent magnitude 5.93 to 6.73 over 7.5 days. It is a yellow-white supergiant ranging between spectral types F7Ib and G2Ib, located around 2037 light-years away. S Muscae is likewise a classical Cepheid, a yellow-white supergiant ranging between spectral types F6Ib and G0Ib and magnitudes 5.89 to 6.49 over a period of 9.66 days. A luminous star around 5.9 times as massive as the Sun, it is a binary star with a blue-white main-sequence star companion likely to be of spectral type B3V to B5V with a mass of just over 5 solar masses, one of the hottest and brightest companions of a cepheid known. The two stars orbit each other every 505 days.
Theta Muscae is a triple star system thought to be around 7,500 light-years distant. It consists of a spectroscopic binary system composed of the Wolf–Rayet star (spectral type: WC5 or 6) and an O-type main-sequence star (spectral type: O6 or O7) that orbit each other every 19 days and a blue supergiant (spectral type: O9.5/B0Iab) set about 46 milliarcseconds apart from them. If the system's estimated distance from Earth is accurate, the binary stars are about 0.5 astronomical units (AU) apart and the supergiant about 100 AU apart from them. All three are highly luminous; combined, they are likely to be over a million times as luminous as the Sun. TU Muscae is a binary star system located around 15,500 light-years away made up of two hot, luminous, blue main-sequence stars of spectral types O7.5V and O9.5V, with masses 23 and 15 times that of the Sun. The stars are so close that they are in contact with each other (overcontact binary) and are classed as a Beta Lyrae variable as their light varies from Earth as they eclipse each other. The system ranges from apparent magnitude 8.17 to 8.75 over around 1.4 days.
Also known as Nova Muscae 1983, GQ Muscae is a binary system consisting of a white dwarf and small star that is about 10% as massive as the Sun. The two orbit each other every 1.4 hours. The white dwarf accumulates material from its companion star via its accretion disc. After a certain amount has accumulated, the star erupts, as it did in 1983, reaching a magnitude of 7.2. Discovered with a magnitude of 7.1 on 18 January 1983, it was the first nova from which X-rays were detected. The soft X-ray transient GRS 1124-683 (also known as Nova Muscae 1991) is a binary object consisting of an orange main-sequence star (GU Muscae) of spectral type K3V–K4V and a black hole of around six solar masses. During the 1991 outburst which led to its discovery, radiation was produced through a process of positron annihilation. GR Muscae is an X-ray source composed of a neutron star of between 1.2 and 1.8 times the mass of the Sun and a low-mass star likely to be around the mass of the Sun in close orbit. Finally, SY Muscae is a symbiotic star system composed of a red giant and white dwarf, where although the larger star is transferring mass to the smaller, no periodic eruption occurs nor does an accretion disc form. The star system varies in magnitude from 10.2 to 12.7 over a period of 624.5 days. V415 Muscae is a nova that had an outburst in 8 June 2022 with an apparent magnitude of +8.7.
Three star systems have been discovered to have exoplanets. HD 111232 is a yellow main-sequence star around 78% as massive as the Sun around 95 light-years distant. It has a planet (HD 111232 b) around 6.8 times the mass of Jupiter that has an orbital period around 1143 days. HD 112410 is a yellow giant of spectral type G8III located around 439 light-years distant. With around 1.54 times the mass of the Sun, it is cooling and expanding along the red-giant branch, having left the main sequence after exhausting its core supply of hydrogen fuel. It has a substellar companion calculated to have a mass 9.2 times that of Jupiter and an orbital period of 124.6 days at a distance around 0.57 AU. Yet another member of the Lower Centaurus Crux subgroup, HD 100546 is a young, blue-white Herbig Ae/Be star of spectral type B9V that has yet to settle on the main sequence—the closest of these stars to Earth around 320 light-years distant. It is surrounded by a circumstellar debris disk from a distance of 0.2 to 4 AU, and again from 13 AU out to a few hundred AU, with evidence for a protoplanet forming at a distance around 47 AU. A gap exists between 4 and 13 AU, which appears to contain a large planet around 20 times the mass of Jupiter, although further examination of the disk profile indicates it might be a more massive object such as a brown dwarf or more than one planet. LP 145-141 is a white dwarf located 15 light-years distant—the fourth-closest to the Solar System. It is considered a good candidate to look for Jupiter-like planets, on account if its proximity and mass.
Deep-sky objects
Located on the border with Circinus is the unusual planetary nebula NGC 5189, estimated to be around 1750 light-years away from Earth. Its complex structure is due to multiple ejections of material from the ageing central star, which are distorted by the presence of a likely binary companion. Located 2.4° east of Eta Muscae is the magnitude-12.9 Engraved Hourglass Nebula (MyCn 18), which lies about 8000 light-years distant from Earth. To Eta's west lies IC 4191, a compact bluish planetary nebula of magnitude 10.6, thought to lie around 10,750 light-years away from Earth. West of Epsilon Muscae is NGC 4071, a large, diffuse planetary nebula of magnitude 12.7 with a magnitude 12 central star, thought to lie around 4000 light-years away from Earth. The Coalsack Nebula is a dark nebula located mainly in neighbouring Crux that intrudes into Musca. NGC 4463 is an open cluster located on its southwestern border. Around five light-years across, it is located around 3400 light-years away.
The comparatively old globular cluster NGC 4833 near Delta Muscae was catalogued by Lacaille in 1755. It is 21,200 light-years distant and somewhat obscured by dust clouds near the galactic plane. The globular cluster NGC 4372 near Gamma Muscae is fainter and likewise partially obscured by dust, but spans more arc minutes. It is 18,900 light-years away from Earth and 23,000 light-years distant from the centre of the Milky Way. Its extremely low metallicity indicates it is very old—one of the oldest clusters in the Milky Way. Extending south from it is the Dark Doodad Nebula, resembling a dark L-shaped river through a bright field of stars. Another dark nebula in the constellation is BHR 71.
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211826 | https://en.wikipedia.org/wiki/Leo%20Minor | Leo Minor | Leo Minor is a small and faint constellation in the northern celestial hemisphere. Its name is Latin for "the smaller lion", in contrast to Leo, the larger lion. It lies between the larger and more recognizable Ursa Major to the north and Leo to the south. Leo Minor was not regarded as a separate constellation by classical astronomers; it was designated by Johannes Hevelius in 1687.
There are 37 stars brighter than apparent magnitude 6.5 in the constellation; three are brighter than magnitude 4.5. 46 Leonis Minoris, an orange giant of magnitude 3.8, is located some 95 light-years from Earth. At magnitude 4.4, Beta Leonis Minoris is the second-brightest star and the only one in the constellation with a Bayer designation. It is a binary star, the brighter component of which is an orange giant and the fainter a yellow-white main sequence star. The third-brightest star is 21 Leonis Minoris, a rapidly rotating white main-sequence star of average magnitude 4.5. The constellation also includes two stars with planetary systems, two pairs of interacting galaxies, and Hanny's Voorwerp, a unique deep-sky object.
History
The classical astronomers Aratus and Ptolemy had noted the region of what is now Leo Minor to be undefined and not containing any distinctive patterns; Ptolemy classified the stars in this area as amorphōtoi (not belonging to a constellation outline) within the constellation Leo.
Johannes Hevelius first depicted Leo Minor in 1687 when he outlined ten new constellations in his star atlas Firmamentum Sobiescianum, and included 18 of its objects in the accompanying Catalogus Stellarum Fixarum. Hevelius decided upon Leo Minor or Leo Junior as a depiction that would align with its beastly neighbours the Lion and the Great Bear. In 1845, English astronomer Francis Baily revised the catalogue of Hevelius's new constellations, and assigned a Greek letter known as Bayer designation to stars brighter than apparent magnitude 4.5. Richard A. Proctor gave the constellation the name Leaena "the Lioness" in 1870, explaining that he sought to shorten the constellation names to make them more manageable on celestial charts.
German astronomer Christian Ludwig Ideler posited that the stars of Leo Minor had been termed Al Thibā' wa-Aulāduhā "Gazelle with her Young" on a 13th-century Arabic celestial globe, recovered by Cardinal Stefano Borgia and housed in the prelate's museum at Velletri. Arabist Friedrich Wilhelm Lach describes a different view, noting that they had been seen as Al Haud "the Pond", which the Gazelle jumps into. In Chinese astronomy, the stars Beta, 30, 37 and 46 Leonis Minoris made up Neiping, a "Court of Judge or Mediator", or Shi "Court Eunuch" or were combined with stars of the neighbouring Leo to make up a large celestial dragon or State Chariot. A line of four stars was known as Shaowei; it represented four Imperial advisors and may have been located in Leo Minor, Leo or adjacent regions.
Characteristics
A dark area of the sky with a triangle of brighter stars just visible to the naked eye in good conditions, Leo Minor has been described by Patrick Moore as having "dubious claims to a separate identity". It is a small constellation bordered by Ursa Major to the north, Lynx to the west, Leo to the south, and touching the corner of Cancer to the southwest. The three-letter abbreviation for the constellation, as adopted by the International Astronomical Union in 1922, is "LMi". The official constellation boundaries, as set by Belgian astronomer Eugène Delporte in 1930, are defined by a polygon of 16 sides. In the equatorial coordinate system, the right ascension coordinates of these borders lie between and , while the declination coordinates are between 22.84° and 41.43°. Ranked 64th out of 88 constellations in size, Leo Minor covers an area of 232.0 square degrees, or 0.562 per cent of the sky. It culminates each year at midnight on 24 February, and at 9 p.m. on 24 May.
Notable features
Stars
There are only three stars in the constellation brighter than magnitude 4.5, and 37 stars with a magnitude brighter than 6.5. Leo Minor does not have a star designated Alpha because Baily erred and allocated a Greek letter to only one star, Beta. It is unclear whether he intended to give 46 Leonis Minoris a Bayer designation, as he recognized Beta and 46 Leonis Minoris as of the appropriate brightness in his catalogue. He died before revising his proofs, which might explain this star's omission.
At magnitude 3.8, the brightest star in Leo Minor is an orange giant of spectral class K0III named 46 Leonis Minoris or Praecipua; its colour is evident when seen through binoculars. Situated from Earth, it has around 32 times the luminosity and is 8.5 times the size of the Sun. It was also catalogued and named as o Leonis Minoris by Johann Elert Bode, which has been misinterpreted as Omicron Leonis Minoris. More confusion occurred with its proper name Praecipua, which appears to have been originally applied to 37 Leonis Minoris in the 1814 Palermo Catalogue of Giuseppe Piazzi, who mistakenly assessed the latter star as the brighter. This name was later connected by Allen with 46 Leonis Minoris—an error perpetuated by subsequent astronomers. The original "Praecipua", 37 Leonis Minoris, has an apparent magnitude of 4.69, but is a distant yellow supergiant of spectral type G2.5IIa and absolute magnitude of −1.84, around distant.
Beta Leonis Minoris is a binary star system. The primary is a giant star of spectral class G9III and apparent magnitude of 4.4. It has around double the mass, 7.8 times the radius and 36 times the luminosity of the Earth's Sun. Separated by 0.1 to 0.6 second of arc from the primary, the secondary is a yellow-white main sequence star of spectral type F8. The two orbit around a common centre of gravity every 38.62 years, and lie away from the Solar System.
Around away and around 10 times as luminous as the Sun, 21 Leonis Minoris is a rapidly rotating white main-sequence star, spinning on its axis in less than 12 hours and very likely flattened in shape. Of average apparent magnitude 4.5 and spectral type A7V, it is a Delta Scuti variable. These are short period (six hours at most) pulsating stars which have been used as standard candles and as subjects to study asteroseismology.
Also known as SU and SV Leonis Minoris, 10 and 11 Leonis Minoris are yellow giants of spectral type G8III, with average magnitudes 4.54 and 5.34 respectively. Both are RS Canum Venaticorum variables, with 10 Leonis Minoris varying by 0.012 magnitude over 40.4 days, and 11 Leonis Minoris by 0.033 magnitude over 18 days. 11 Leonis Minoris has a red dwarf companion of spectral type M5V and apparent magnitude 13.0. 20 Leonis Minoris is a multiple star system away from the Sun. The main star is another yellow star, this time a dwarf of spectral type G3Va and apparent magnitude 5.4. The companion is an old, active red dwarf that has a relatively high metallicity and is of spectral type M6.5. The fact that the secondary star is brighter than expected indicates it is likely two stars very close together that are unable to be made out separately with current viewing technology.
R and S Leonis Minoris are long-period Mira variables, while U Leonis Minoris is a semiregular variable; all three are red giants of spectral types M6.5e-M9.0e, M5e and M6 respectively. R varies between magnitudes 6.3 and 13.2 during a period of 372 days, S varies between magnitudes 8.6 and 13.9 during a period of 234 days, and U varies between magnitudes 10.0 and 13.3 during a period of 272 days. The lack of bright stars makes finding these objects challenging for amateur astronomers. G 117-B15A, also known as RY Leonis Minoris, is a pulsating white dwarf of apparent magnitude 15.5. With a period of approximately 215 seconds, and losing a second every 8.9 million years, the 400-million-year-old star has been proposed as the most stable celestial clock.
SX Leonis Minoris is a dwarf nova of the SU Ursae Majoris type that was identified in 1994. It consists of a white dwarf and a donor star, which orbit each other every 97 minutes. The white dwarf sucks matter from the other star onto an accretion disc and heats up to between 6000 and 10000 K. The dwarf star erupts every 34 to 64 days, reaching magnitude 13.4 in these outbursts and remaining at magnitude 16.8 when quiet. Leo Minor contains another dwarf nova, RZ Leonis Minoris, which brightens to magnitude 14.2 from a baseline magnitude of around 17 but does so at shorter intervals than other dwarf novae.
Two stars with planetary systems have been found. HD 87883 is an orange dwarf of magnitude 7.57 and spectral type K0V 18 parsecs distant from Earth. With a diameter three quarters that of Earth's sun, it is only 31 per cent as luminous. It is orbited by a planet around 1.78 times the mass of Jupiter every 7.9 years, and there are possibly other smaller planets. HD 82886 (Illyrian) is a yellow dwarf of spectral type G0 and visual magnitude 7.63. A planet 1.3 times the mass of Jupiter and orbiting every 705 days was discovered in 2011.
Deep-sky objects
In terms of deep-sky objects, Leo Minor contains many galaxies viewable in amateur telescopes. Located 3 degrees southeast of 38 Leonis Minoris, NGC 3432 is seen nearly edge on. It is of apparent magnitude 11.7 and measures 6.8 by 1.4 arcminutes. Located 42 million light-years away, it is moving away from the Solar System at a rate of 616 km per second. In 2000, a star within the galaxy brightened to magnitude 17.4, and has since been determined to be a luminous blue variable and supernova impostor. It has tidal filaments and intense star formation, so it was listed in Halton Arp's Atlas of Peculiar Galaxies. NGC 3003, a SBbc barred spiral galaxy with an apparent magnitude of 12.3 and an angular size of 5.8 arcminutes, is seen almost edge-on. NGC 3344, 25 million light-years distant, is face-on towards Earth. Measuring 7.1 by 6.5 arcminutes in size, it has an apparent magnitude of 10.45. NGC 3504 is a starburst barred spiral galaxy of apparent magnitude 11.67 and measuring 2.1 by 2.7 arcminutes. It has hosted supernovae in 1998 and 2001. It and the spiral galaxy NGC 3486 are also almost face-on towards Earth; the latter is of magnitude 11.05 and measures 7.1 by 5.2 arcminutes. NGC 2859 is an SB0-type lenticular galaxy.
At least two pairs of interacting galaxies have been observed. Arp 107 is a pair of galaxies in the process of merging, located 450 million light-years away. NGC 3395 and NGC 3396 are a spiral and irregular barred spiral galaxy, respectively, that are interacting, located 1.33 degrees southwest of 46 Leonis Minoris.
The unique deep-sky object known as Hanny's Voorwerp was discovered in Leo Minor in 2007 by Dutch school teacher Hanny van Arkel while participating as a volunteer in the Galaxy Zoo project. Lying near the 650-million-light-year-distant spiral galaxy IC 2497, it is around the same size as the Milky Way. It contains a 16,000-light-year-wide hole. The voorwerp is thought to be the visual light echo of a quasar now gone inactive, possibly as recently as 200,000 years ago.
Meteor shower
Discovered by Dick McCloskey and Annette Posen of the Harvard Meteor Program in 1959, the Leonis Minorid meteor shower peaks between 18 and 29 October. The shower's parent body is the long period comet C/1739 K1 (Zanotti). It is a minor shower, and can only be seen from the Northern Hemisphere.
In Popular Culture
The fact that Leo Minor was not a named constellation until 1687 provided proof that an alleged play by Euripides was in fact a forgery, in the Inspector Lewis episode "The Lions of Nemea".
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211828 | https://en.wikipedia.org/wiki/Columba%20%28constellation%29 | Columba (constellation) | Columba is a faint constellation designated in the late sixteenth century, remaining in official use, with its rigid limits set in the 20th century. Its name is Latin for dove. It takes up 1.31% of the southern celestial hemisphere and is just south of Canis Major and Lepus.
History
Early 3rd century BC: Aratus's astronomical poem Phainomena (lines 367–370 and 384–385) mentions faint stars where Columba is now but does not fit any name or figure to them.
2nd century AD: Ptolemy lists 48 constellations in the Almagest. While Columba is not yet among them, several stars south of Canis Major listed in this work will eventually become part of Columba.
c. 150–215 AD: Clement of Alexandria wrote in his Logos Paidogogos"Αἱ δὲ σφραγῖδες ἡμῖν ἔστων πελειὰς ἢ ἰχθὺς ἢ ναῦς οὐριοδρομοῦσα ἢ λύρα μουσική, ᾗ κέχρηται Πολυκράτης, ἢ ἄγκυρα ναυτική," (= "[when recommending symbols for Christians to use], let our seals be a dove or a fish or a ship running in a good wind or a musical lyre ... or a ship's anchor ..."), with no mention of stars or astronomy.
1592 AD: Petrus Plancius first depicted Columba on the small celestial planispheres of his large wall map to differentiate the 'unformed stars' of the large constellation Canis Major. Columba is also shown on his smaller world map of 1594 and on early Dutch celestial globes. Plancius named the constellation Columba Noachi ("Noah's Dove"), referring to the dove that gave Noah the information that the Great Flood was receding. This name is found on early 17th-century celestial globes and star atlases.
1603: Frederick de Houtman listed Columba as "De Duyve med den Olijftack" ("the dove with the olive branch")
1603: Bayer's sky atlas Uranometria was published. It includes Columba as Columba Noachi.
1624: Bartschius listed Columba in his Usus Astronomicus as "Columba Nohae".
1662: Caesius published Coelum Astronomico-Poeticum, including an inaccurate Latin translation of the above text of Clement of Alexandria: it mistranslated "ναῦς οὐριοδρομοῦσα" as Latin "Navis coelestis cursu in coelum tendens" ("Ship of the sky following a course in the sky"), perhaps misunderstanding "οὐριο-" as "up in the air or sky" by analogy with οὐρανός = "sky".
1679: Halley mentioned Columba in his work Catalogus Stellarum Australium from his observations on St. Helena.
1679: Augustin Royer published a star atlas that showed Columba as a constellation.
c.1690: Hevelius's Prodromus Astronomiae showed Columba but did not list it as a constellation.
1712 (pirated) and 1725 (authorized): Flamsteed's work Historia Coelestis Britannica showed Columba but did not list it as a constellation.
1757 or 1763: Lacaille listed Columba as a constellation and catalogued its stars.
1889: Richard H. Allen, misled by Caesius's mistranslation, wrote that the Columba asterism may have been invented in Roman/Greek times, but with a footnote saying that it may have been another star group.
2019: OSIRIS-REx students discovered a black hole in the constellation Columba, based on observing X-ray bursts.
Features
Stars
Columba is rather inconspicuous with the brightest star, Alpha Columbae, being only of magnitude 2.7. This, a blue-white star, has a pre-Bayer, traditional, Arabic name Phact (meaning ring dove) and is 268 light-years from Earth. The only other named star is Beta Columbae, which has the alike-status name Wazn. It is an orange-hued giant star of magnitude 3.1, 87 light-years away. The constellation contains the runaway star μ Columbae. Exoplanet NGTS-1b and its star NGTS-1 are in Columba.
General radial velocity
Columba contains the solar antapex – the opposite to the net direction of the solar system
Deep-sky objects
The globular cluster NGC 1851 appears in Columba at 7th magnitude in a far part of our galaxy at 39,000 light-years away - it is resolvable south of at greatest latitude +40°N in medium-sized amateur telescopes (under good conditions).
| Physical sciences | Other | Astronomy |
211830 | https://en.wikipedia.org/wiki/Horologium%20%28constellation%29 | Horologium (constellation) | Horologium (Latin , the pendulum clock, from Greek , ) is a constellation of six stars faintly visible in the southern celestial hemisphere. It was first described by the French astronomer Nicolas-Louis de Lacaille in 1756 and visualized by him as a clock with a pendulum and a second hand. In 1922 the constellation was redefined by the International Astronomical Union (IAU) as a region of the celestial sphere containing Lacaille's stars, and has since been an IAU designated constellation. Horologium's associated region is wholly visible to observers south of 23°N.
The constellation's brightest star—and the only one brighter than an apparent magnitude of 4—is Alpha Horologii (at 3.85), an aging orange giant star that has swollen to around 11 times the diameter of the Sun. The long-period variable-brightness star, R Horologii (4.7 to 14.3), has one of the largest variations in brightness among all stars in the night sky visible to the unaided eye. Four star systems in the constellation are known to have exoplanets; at least one—Gliese 1061—contains an exoplanet in its habitable zone.
History
The French astronomer Nicolas-Louis de Lacaille first described the constellation as l'Horloge à pendule & à secondes (Clock with pendulum and seconds hand) in 1756, after he had observed and catalogued almost 10,000 southern stars during a two-year stay at the Cape of Good Hope. He devised fourteen new constellations in previously uncharted regions of the southern celestial hemisphere, which were not visible from Europe. All but one honoured scientific instruments, and so symbolised the Age of Enlightenment. The constellation name was Latinised to Horologium in a catalogue and updated chart published posthumously in 1763. The Latin term is ultimately derived from the Ancient Greek ὡρολόγιον, for an instrument for telling the hour.
Characteristics
Covering a total of 248.9 square degrees or 0.603% of the sky, Horologium ranks 58th in area out of the 88 modern constellations. Its position in the southern celestial hemisphere means the whole constellation is visible to observers south of 23°N. Horologium is bordered by five constellations: Eridanus (the Po river or Nile river), Caelum (the chisel), Reticulum (the reticle), Dorado (the dolphin/swordfish), and Hydrus (the male water snake). The three letter abbreviation for the constellation, as adopted by the International Astronomical Union in 1922, is "Hor". The official constellation boundaries are defined by a twenty-two-sided polygon (illustrated in infobox). In the equatorial coordinate system, the right ascension coordinates of these borders lie between and , while the declination coordinates are between −39.64° and −67.04°.
Features
Stars
Horologium has one star brighter than apparent magnitude 4, and 41 stars brighter than or equal to magnitude 6.5. Lacaille charted and designated 11 stars in the constellation, giving them the Bayer designations Alpha (α Hor) through Lambda Horologii (λ Hor) in 1756. In the mid-19th century, English astronomer Francis Baily removed the designations of two—Epsilon and Theta Horologii—as he held they were too faint to warrant naming. He was unable to find a star that corresponded to the coordinates of Lacaille's Beta Horologii. Determining that the coordinates were wrong, he assigned the designation to another star. Kappa Horologii, too, was unable to be verified—although it most likely was the star HD 18292—and the name fell out of use. In 1879, American astronomer Benjamin Apthorp Gould assigned designations to what became Mu and Nu Horologii as he felt they were bright enough to warrant them.
At magnitude 3.9, Alpha Horologii is the brightest star in the constellation, located 115 (±0.5) light-years from Earth. German astronomer Johann Elert Bode depicted it as the pendulum of the clock, while Lacaille made it one of the weights. It is an orange giant star of spectral type K2III that has swollen to around 11 times the diameter of the Sun, having spent much of its life as a white main-sequence star. At an estimated 1.55 times the mass of the Sun, it is radiating 38 times the Sun's luminosity from its photosphere at an effective (surface) temperature of 5,028K.
At magnitude 4.93, Delta Horologii is the second-brightest star in the constellation, and forms a wide optical double with Alpha. Delta itself is a true binary system composed of a white main sequence star of spectral type A5V that is 1.41 times as massive as the Sun with a magnitude of 5.15 and its fainter companion of magnitude 7.29. The system is located 179 (±4) light-years from the earth.
At magnitude 5.0, Beta Horologii is a white giant 63 times as luminous as the Sun with an effective temperature of 8,303K. It is 312 (±4) light-years from Earth, and has been little-studied. Lambda Horologii is an ageing yellow-white giant star of spectral type F2III that spins around at 140km/second, and is hence mildly flattened at its poles (oblate). It is 161 (±1) light-years from Earth.
With a magnitude of 5.24, Nu Horologii is a white main sequence star of spectral type A2V located 169 (±1) light-years from Earth that is around 1.9 times as massive as the Sun. Estimated to be around 540 million years old, it has a debris disk that appears to have two components: an inner disk is orbiting at a distance of , while an outer disk lies from the star. The estimated mass of the disks is the mass of the Earth.
Horologium has several variable stars. R Horologii is a red giant Mira variable with one of the widest ranges in brightness known of stars in the night sky visible to the unaided eye. It is around 1,000 light-years from Earth. It has a minimum magnitude of 14.3 and a maximum magnitude of 4.7, with a period of approximately 13 months. T and U Horologii are also Mira variables. The Astronomical Society of Southern Africa reported in 2003 that observations of these two stars were needed as data on their light curves was incomplete. TW Horologii is a semiregular variable red giant star that is classified as a carbon star, and is 1,370 (±70) light-years from Earth.
Iota Horologii is a yellow-white dwarf star 1.23 (±0.12) times as massive and 1.16 (±0.04) times as wide as the Sun with a spectral type of F8V, 57 (±0.05) light-years from Earth. Its chemical profile, movement and age indicate it formed within the Hyades cluster but has drifted around 130 light-years away from the other members. It has a planet at least 2.5 times as massive as Jupiter orbiting it every 307 days. HD 27631 is a Sun-like star located 164 (±0.3) light-years from Earth which was found to have a planet at least 1.45 times as massive as Jupiter that takes 2,208 (±66) days (six years) to complete an orbit. WASP-120 is a yellow-white main-sequence star around 1.4 times as massive as the Sun with a spectral type of F5V that is estimated to be 2.6 (±0.5) billion years old. It has a massive planet around 4.85 times the mass of Jupiter that completes its orbit every 3.6 days, and has an estimated surface temperature of 1,880 (±70)K.
With an apparent magnitude of 13.06, Gliese 1061 is a red dwarf of spectral type M5.5V that has 12% of the mass and 15% of the diameter of the Sun, and shines with only 0.17% of its luminosity. Located 12 light-years away from Earth, it is the 20th-closest single star or stellar system to the Sun. In August 2019, it was announced that it had three planets, one of which lay in its habitable zone.
Deep-sky objects
Horologium is home to many deep-sky objects, including several globular clusters. NGC 1261 is a globular cluster of magnitude 8, located 53,000 light-years from Earth. It lies 4.7 degrees north-northeast of Mu Horologii. The globular cluster Arp-Madore 1 is the most remote known globular cluster in the Milky Way at a distance of from Earth.
NGC 1512 is a barred spiral galaxy 2.1 degrees west-southwest of Alpha Horologii with an apparent magnitude of 10.2. About five arcmin (13.8 kpc) away is the dwarf lenticular galaxy NGC 1510. The two are in the process of a merger which has been going on for 400 million years.
The Horologium-Reticulum Supercluster is a galaxy supercluster, second in size only to the Shapley Supercluster in the local universe (anything within 200 mpc of Earth). It contains over 20 Abell galaxy clusters and covers more than 100 deg2 of the sky, centered roughly at equatorial coordinates α = , δ = .
| Physical sciences | Other | Astronomy |
211833 | https://en.wikipedia.org/wiki/Indus%20%28constellation%29 | Indus (constellation) | Indus is a constellation in the southern sky first professionally surveyed by Europeans in the 1590s and mapped on a globe by Petrus Plancius by early 1598. It was included on a plate illustrating southern constellations in Bayer's sky atlas Uranometria in 1603. It lies well south of the Tropic of Capricorn but its triangular shape can be seen for most of the year from the Equator. It is elongated from north to south and has a complex boundary. The English translation of its name is generally given as the Indian, though it is unclear which indigenous people the constellation was originally supposed to represent.
Features
Indus lacks stars of the top 100 in brightness viewed from the solar system (apparent magnitude). Two of its stars rank of third magnitude and three of fourth magnitude.
Bright stars
Alpha Indi, its brightest, is an orange giant of magnitude 3.1, 101 light-years away. Beta Indi is an orange giant of magnitude 3.7, 600 light-years distant. Delta Indi is a white star of magnitude 4.4, 185 light-years from Earth. The three form a near-perfect right-angled triangle, such that Beta marks the right angle and is in the south-east.
Epsilon Indi is one of the closest stars to Earth, approximately 11.8 light years away. It is an orange dwarf of magnitude 4.7, meaning that the yellow dwarf Sun is slightly hotter and larger. The system has been discovered to contain a pair of binary brown dwarfs, and has long been a prime candidate in SETI studies. This star has the third-highest proper motion of all visible to the unaided eye, as ranks behind Groombridge 1830 and 61 Cygni, and the ninth-highest overall. This will move the star into Tucana around 2640. It figures directly between Alpha and Beta.
Indus is home to one bright binary star. Theta Indi is a binary star divisible in small amateur telescopes, 97 light-years from Earth. Its primary is a white star of magnitude 4.5 and its secondary is a white star of magnitude 7.0. It figures close to the hypotenuse of the right-angled triangle of Alpha, Beta and Delta, the three brightest stars of Indus.
Variable stars
T Indi is the only bright variable star in Indus. It is a semi-regular, deeply coloured red giant with a period of 11 months, 1900 light-years away. Its minimum magnitude is 7 and its maximum: 5.
Galaxies
Galaxies include NGC 7038, NGC 7049, and NGC 7090.
Supernovae
All Sky Automated Survey for SuperNovae (ASAS-SN) in 2015 detected a superluminous supernova, named ASASSN-15lh (also designated SN 2015L). Based on the study conducted by Subo Dong and team from the Kavli Institute for Astronomy and Astrophysics (KIAA) at Peking University, it was approximately doubly luminous to any supernova detected, and at peak was almost 50 times more intrinsically luminous than the Milky Way.
Its distance: approximately 3.82 gigalight-years, denoting an age approximately half that of the universe.
History
The constellation was created by Petrus Plancius who made a fairly large celestial globe from the observations of Pieter Dirkszoon Keyser and Frederick de Houtman. The first depiction of this constellation in a celestial atlas followed in Johann Bayer's Uranometria of 1603. Plancius portrayed the figure as a nude male with three arrows in one hand and one in the other, as a native, lacking quiver and bow. It is among the twelve constellations introduced by Keyser and de Houtman, which first appeared on a celestial globe in 1598.
| Physical sciences | Other | Astronomy |
211857 | https://en.wikipedia.org/wiki/Banana%20slug | Banana slug | Banana slugs (Ariolimax) are a genus of air-breathing, terrestrial gastropod slugs in the family Ariolimacidae. Their often yellow color and elongated shape can resemble a banana, thus giving rise to their common name.
Species
Species within the genus Ariolimax include:
Synonyms:
Ariolimax andersonii J. G. Cooper, 1872: synonym of Prophysaon andersonii (J.G. Cooper, 1872) (unaccepted combination)
Ariolimax californicus J.G. Cooper, 1872: synonym of Ariolimax californicus J.G. Cooper, 1872
Ariolimax californicus costaricensis Cockerell, 1890: synonym of Deroceras costaricensis (Cockerell, 1890)
Ariolimax hemphilli W. G. Binney, 1875: synonym of Hesperarion hemphilli (W.G. Binney, 1875)
Ariolimax niger J. G. Cooper, 1872: synonym of Hesperarion niger (J. G. Cooper, 1872)
Ariolimax steindachneri Babor, 1900: synonym of Ariolimax columbianus (Gould, 1851)
Description
Banana slugs are often bright yellow (giving rise to the banana nomenclature) although they may also be greenish, brown, tan, or white. The species Ariolimax columbianus sometimes has black spots that are so extensive that the animal looks almost entirely black. Individual slugs can change colors with alterations in diet, light exposure, and moisture levels. Color may also be an indicator of the age or health status of an individual.
The Pacific banana slug is the second-largest species of terrestrial slug in the world, achieving a length of up to and a weight of up to . The largest slug species is Limax cinereoniger, which can reach in length. Banana slugs have an average lifespan of 1–7 years.
Banana slugs (like other gastropods and many other mollusks) have a radula, a ribbon-like anatomical structure covered in rows of microscopic teeth. The radula is used for feeding. Individuals can move at per minute.
Slugs use two pairs of tentacles to sense their environment. The longer, upper pair are used to detect light or movement. The shorter, lower pair are used to detect chemicals. Both pairs of tentacles are muscular hydrostats which can be fully retracted when necessary to avoid damage. If the animal loses a tentacle, the slug can grow a new one.
Banana slugs have a single pallial lung which opens externally via a pneumostome on the right side of the mantle of the animal. The pneumostome lung cavity is heavily vascularized to allow gas exchange. Dehydration is a major problem for the mollusk; to combat this, banana slugs excrete a thick coating of mucus around their bodies and can also aestivate. To do so, they secrete a protective layer of mucus and insulate themselves with a layer of soil and leaves. They remain inactive in this state until the environment becomes moist again. Due to their susceptibility to desiccation, they are more commonly active at night, but also appear during cool, moist days.
The slime also contains pheromones to attract other slugs for mating. Slugs are simultaneous hermaphrodites, and reproduce by exchanging sperm with their mate. They produce up to 75 translucent eggs, which are laid in a log or on leaves. Slugs mate and lay eggs throughout the year. The adults provide no further care for their eggs beyond finding a suitable hiding spot, and the eggs are abandoned as soon as the clutch is laid.
Distribution
Ariolimax columbianus is native to the forest floors along North America's Pacific coastal coniferous rainforest belt (including douglas-fir forests and redwood forests) which stretches from Southeastern Alaska to Santa Cruz, California.
Several discontinuous populations also occur in forested slopes of the coastal and transverse mountain ranges south of Santa Cruz as far south as Ventura County, with a tiny, isolated population located in Palomar Mountain State Park within the Palomar Mountain Range in San Diego County, California. The Palomar Mountains have lush Sierra Nevada-like coniferous forests and black oak woodlands unlike the surrounding semiarid lands of inland San Diego County and mark the southernmost population of banana slugs.
Small, isolated populations also occur east of the Pacific Coast such as in the inland coniferous rainforests of British Columbia's Columbia Mountains (interior wet-belt), just west of the Canadian Rockies, and have been seen at lower elevations near creeks and damp areas of Mount Revelstoke National Park. Small populations of banana slugs have also been seen along creek and damp areas of the western slopes of the Sierra Nevada mountains to the north of Yosemite National Park in California. Slug densities in these outlying areas in the Columbia Mountains, Sierra Nevada Mountains, and areas south of Santa Cruz are low compared to densities in the coastal coniferous rainforest belt and are rather restricted to damp areas near creeks, ravines, and gullies. This population may also be a relic from the Pleistocene epoch.
Ecology
Banana slugs are detritivores, or decomposers. They consume leaves, animal droppings, moss, and dead plant material, and then recycle them into soil humus. They are generalist feeders, though they exhibit a preference for certain mushrooms. Through their consumption of various plant matter and animal excretion, banana slugs acts as agents of seed dispersal. They spread seeds and spores, and excrete a nitrogen-rich fertilizer. Additionally, by consuming dead organic matter, they contribute to decomposition and the nutrient cycles and are an important aspect of the ecosystem.
Small mammals, snakes, and salamanders sometimes eat banana slugs.
Slime
The mucus secreted by banana slugs contains chemicals that can numb the tongue of predators. This mucus can absorb up to 100 times its volume in water. Technically, this slime is neither liquid nor solid, but rather a liquid crystal substance. Slug slime bears a resemblance to human mucus, primarily composed of proteins known as mucins. Upon contact with environmental moisture, dry mucin granules swell dramatically, expanding to hundreds of times their original volume. This expansion contributes to the slimy texture but doesn't involve the numbing effect seen in banana slugs. Moreover, the slime serves the purpose of maintaining moisture on the skin of these slug-like creatures, akin to fruit, which facilitates gas exchange to support respiration. The glycoprotein-based mucous behaves as a liquid crystal, occupying a unique state between a solid and a liquid. Its molecular dispersion prevents it from being a true liquid, yet it lacks the compactness of a solid, displaying an organized, crystal-like pattern.
This mucous secretion serves a multitude of purposes in the life of a banana slug. It aids in locomotion, allowing them to glide gracefully across the forest floor. The slime also serves as a protective barrier, acting as a deterrent against potential predators and pathogens. Additionally, banana slug slime plays a role in reproduction, as slugs deposit chemical cues within the slime to attract potential mates. The slime's vivid yellow color, which is responsible for their common name, not only adds to their eccentric appearance but also acts as a warning signal to potential predators that it may not be a tasty meal. Banana slug mucus has potential to carry mites and other parasitic organisms.
Reproduction
Banana slugs are hermaphroditic, possessing both male and female reproductive organs. This feature allows them to mate as either sex or even fertilize themselves. The versatile slime also plays a role in their mating rituals. During the mating phase, banana slugs release pheromones through their adhesive trails to attract potential partners. Due to a slug's ability for self-fertilization, there is a high degree of homozygosity. Notably, their courtship process can be quite robust and, in certain species, may culminate in apophallation. During the mating ritual between two slugs, they both engage in a dance-like action where they will circle each other, occasionally lunging, nipping, and tail whipping. University of California, Santa Cruz research associate Janet Leonard notes that apophallation occurs some time after copulation, and the process of chewing takes around 45 minutes, but can take longer. After mating, the slugs separate and seek a sheltered location to lay their eggs. An impregnated slug will typically lay a range of 3 to 50 eggs in a small crevice. These eggs will typically hatch within a month or two, and the young slugs will mature over several months. Their offspring will typically contain a pale white hue and as they grow acquire their distinct yellow color.
In culture
Banana slugs have been used as food by Yurok Natives of the North Coast and by German immigrants in the 19th and early 20th centuries. A yearly festival and contest is held at Russian River including slug races and a contest for recipes though, even when fed corn meal to purge them or soaked in vinegar to remove slime, the slugs' flavor is not always well regarded for the modern palate, and the most successful entries are often those in which the flavor is unnoticeable.
The banana slug (specifically Ariolimax dolichophallus, which is the most common banana slug on campus) is the mascot of the University of California, Santa Cruz. It is common in local forests, was approved by a student vote, and was named the best college mascot by Reader's Digest in 2004. The “Every Little Thing” podcast also dedicated half an episode to discussing the banana slug and its historical relevance to the UCSC. After celebrating 25 years of the establishment of the banana slug as UCSC’s mascot, on Sept. 27, 2011, the Santa Cruz City Council declared that day as the official “Day of the UC Santa Cruz Banana Slug.” One T-shirt caption was: "No known predators!"
On September 27, 2024, Governor Gavin Newsom signed AB 1850, making the banana slug the official state slug of California.
| Biology and health sciences | Gastropods | Animals |
211922 | https://en.wikipedia.org/wiki/Impulse%20%28physics%29 | Impulse (physics) | In classical mechanics, impulse (symbolized by or Imp) is the change in momentum of an object. If the initial momentum of an object is , and a subsequent momentum is , the object has received an impulse :
Momentum is a vector quantity, so impulse is also a vector quantity:
Newton’s second law of motion states that the rate of change of momentum of an object is equal to the resultant force acting on the object:
so the impulse delivered by a steady force acting for time is:
The impulse delivered by a varying force is the integral of the force with respect to time:
The SI unit of impulse is the newton second (N⋅s), and the dimensionally equivalent unit of momentum is the kilogram metre per second (kg⋅m/s). The corresponding English engineering unit is the pound-second (lbf⋅s), and in the British Gravitational System, the unit is the slug-foot per second (slug⋅ft/s).
Mathematical derivation in the case of an object of constant mass
Impulse produced from time to is defined to be
where is the resultant force applied from to .
From Newton's second law, force is related to momentum by
Therefore,
where is the change in linear momentum from time to . This is often called the impulse-momentum theorem (analogous to the work-energy theorem).
As a result, an impulse may also be regarded as the change in momentum of an object to which a resultant force is applied. The impulse may be expressed in a simpler form when the mass is constant:
where
is the resultant force applied,
and are times when the impulse begins and ends, respectively,
is the mass of the object,
is the final velocity of the object at the end of the time interval, and
is the initial velocity of the object when the time interval begins.
Impulse has the same units and dimensions as momentum. In the International System of Units, these are . In English engineering units, they are .
The term "impulse" is also used to refer to a fast-acting force or impact. This type of impulse is often idealized so that the change in momentum produced by the force happens with no change in time. This sort of change is a step change, and is not physically possible. However, this is a useful model for computing the effects of ideal collisions (such as in videogame physics engines). Additionally, in rocketry, the term "total impulse" is commonly used and is considered synonymous with the term "impulse".
Variable mass
The application of Newton's second law for variable mass allows impulse and momentum to be used as analysis tools for jet- or rocket-propelled vehicles. In the case of rockets, the impulse imparted can be normalized by unit of propellant expended, to create a performance parameter, specific impulse. This fact can be used to derive the Tsiolkovsky rocket equation, which relates the vehicle's propulsive change in velocity to the engine's specific impulse (or nozzle exhaust velocity) and the vehicle's propellant-mass ratio.
| Physical sciences | Classical mechanics | null |
211941 | https://en.wikipedia.org/wiki/B%20cell | B cell | B cells, also known as B lymphocytes, are a type of white blood cell of the lymphocyte subtype. They function in the humoral immunity component of the adaptive immune system. B cells produce antibody molecules which may be either secreted or inserted into the plasma membrane where they serve as a part of B-cell receptors. When a naïve or memory B cell is activated by an antigen, it proliferates and differentiates into an antibody-secreting effector cell, known as a plasmablast or plasma cell. In addition, B cells present antigens (they are also classified as professional antigen-presenting cells, APCs) and secrete cytokines. In mammals B cells mature in the bone marrow, which is at the core of most bones. In birds, B cells mature in the bursa of Fabricius, a lymphoid organ where they were first discovered by Chang and Glick, which is why the B stands for bursa and not bone marrow, as commonly believed.
B cells, unlike the other two classes of lymphocytes, T cells and natural killer cells, express B cell receptors (BCRs) on their cell membrane. BCRs allow the B cell to bind to a foreign antigen, against which it will initiate an antibody response. B cell receptors are extremely specific, with all BCRs on a B cell recognizing the same epitope.
Development
B cells develop from hematopoietic stem cells (HSCs) that originate from bone marrow. HSCs first differentiate into multipotent progenitor (MPP) cells, then common lymphoid progenitor (CLP) cells. From here, their development into B cells occurs in several stages (shown in image to the right), each marked by various gene expression patterns and immunoglobulin H chain and L chain gene loci arrangements, the latter due to B cells undergoing V(D)J recombination as they develop.
B cells undergo two types of selection while developing in the bone marrow to ensure proper development, both involving B cell receptors (BCR) on the surface of the cell. Positive selection occurs through antigen-independent signalling involving both the pre-BCR and the BCR. If these receptors do not bind to their ligand, B cells do not receive the proper signals and cease to develop. Negative selection occurs through the binding of self-antigen with the BCR; if the BCR can bind strongly to self-antigen, then the B cell undergoes one of four fates: clonal deletion, receptor editing, anergy, or ignorance (B cell ignores signal and continues development). This negative selection process leads to a state of central tolerance, in which the mature B cells do not bind self antigens present in the bone marrow.
To complete development, immature B cells migrate from the bone marrow into the spleen as transitional B cells, passing through two transitional stages: T1 and T2. Throughout their migration to the spleen and after spleen entry, they are considered T1 B cells. Within the spleen, T1 B cells transition to T2 B cells. T2 B cells differentiate into either follicular (FO) B cells or marginal zone (MZ) B cells depending on signals received through the BCR and other receptors. Once differentiated, they are now considered mature B cells, or naïve B cells.
Activation
B cell activation occurs in the secondary lymphoid organs (SLOs), such as the spleen and lymph nodes. After B cells mature in the bone marrow, they migrate through the blood to SLOs, which receive a constant supply of antigen through circulating lymph. At the SLO, B cell activation begins when the B cell binds to an antigen via its BCR. Although the events taking place immediately after activation have yet to be completely determined, it is believed that B cells are activated in accordance with the kinetic segregation model , initially determined in T lymphocytes. This model denotes that before antigen stimulation, receptors diffuse through the membrane coming into contact with Lck and CD45 in equal frequency, rendering a net equilibrium of phosphorylation and non-phosphorylation. It is only when the cell comes in contact with an antigen presenting cell that the larger CD45 is displaced due to the close distance between the two membranes. This allows for net phosphorylation of the BCR and the initiation of the signal transduction pathway. Of the three B cell subsets, FO B cells preferentially undergo T cell-dependent activation while MZ B cells and B1 B cells preferentially undergo T cell-independent activation.
B cell activation is enhanced through the activity of CD21, a surface receptor in complex with surface proteins CD19 and CD81 (all three are collectively known as the B cell coreceptor complex). When a BCR binds an antigen tagged with a fragment of the C3 complement protein, CD21 binds the C3 fragment, co-ligates with the bound BCR, and signals are transduced through CD19 and CD81 to lower the activation threshold of the cell.
T cell-dependent activation
Antigens that activate B cells with the help of T-cell are known as T cell-dependent (TD) antigens and include foreign proteins. They are named as such because they are unable to induce a humoral response in organisms that lack T cells. B cell responses to these antigens takes multiple days, though antibodies generated have a higher affinity and are more functionally versatile than those generated from T cell-independent activation.
Once a BCR binds a TD antigen, the antigen is taken up into the B cell through receptor-mediated endocytosis, degraded, and presented to T cells as peptide pieces in complex with MHC-II molecules on the cell membrane. T helper (TH) cells, typically follicular T helper (TFH) cells recognize and bind these MHC-II-peptide complexes through their T cell receptor (TCR). Following TCR-MHC-II-peptide binding, T cells express the surface protein CD40L as well as cytokines such as IL-4 and IL-21. CD40L serves as a necessary co-stimulatory factor for B cell activation by binding the B cell surface receptor CD40, which promotes B cell proliferation, immunoglobulin class switching, and somatic hypermutation as well as sustains T cell growth and differentiation. T cell-derived cytokines bound by B cell cytokine receptors also promote B cell proliferation, immunoglobulin class switching, and somatic hypermutation as well as guide differentiation. After B cells receive these signals, they are considered activated.
Once activated, B cells participate in a two-step differentiation process that yields both short-lived plasmablasts for immediate protection and long-lived plasma cells and memory B cells for persistent protection. The first step, known as the extrafollicular response, occurs outside lymphoid follicles but still in the SLO. During this step activated B cells proliferate, may undergo immunoglobulin class switching, and differentiate into plasmablasts that produce early, weak antibodies mostly of class IgM.
The second step consists of activated B cells entering a lymphoid follicle and forming a germinal center (GC), which is a specialized microenvironment where B cells undergo extensive proliferation, immunoglobulin class switching, and affinity maturation directed by somatic hypermutation. These processes are facilitated by TFH and follicular dendritic cells within the GC and generate both high-affinity memory B cells and long-lived plasma cells. Resultant plasma cells secrete large numbers of antibodies and either stay within the SLO or, more preferentially, migrate to bone marrow.
T cell-independent activation
Antigens that activate B cells without T cell help are known as T cell-independent (TI) antigens and include foreign polysaccharides and unmethylated CpG DNA. They are named as such because they are able to induce a humoral response in organisms that lack T cells. B cell response to these antigens is rapid, though antibodies generated tend to have lower affinity and are less functionally versatile than those generated from T cell-dependent activation.
As with TD antigens, B cells activated by TI antigens need additional signals to complete activation, but instead of receiving them from T cells, they are provided either by recognition and binding of a common microbial constituent to toll-like receptors (TLRs) or by extensive crosslinking of BCRs to repeated epitopes on a bacterial cell. B cells activated by TI antigens go on to proliferate outside lymphoid follicles but still in SLOs (GCs do not form), possibly undergo immunoglobulin class switching, and differentiate into short-lived plasmablasts that produce early, weak antibodies mostly of class IgM, but also some populations of long-lived plasma cells.
Memory B cell activation
Memory B cell activation begins with the detection and binding of their target antigen, which is shared by their parent B cell. Some memory B cells can be activated without T cell help, such as certain virus-specific memory B cells, but others need T cell help. Upon antigen binding, the memory B cell takes up the antigen through receptor-mediated endocytosis, degrades it, and presents it to T cells as peptide pieces in complex with MHC-II molecules on the cell membrane. Memory T helper (TH) cells, typically memory follicular T helper (TFH) cells, that were derived from T cells activated with the same antigen recognize and bind these MHC-II-peptide complexes through their TCR. Following TCR-MHC-II-peptide binding and the relay of other signals from the memory TFH cell, the memory B cell is activated and differentiates either into plasmablasts and plasma cells via an extrafollicular response or enter a germinal center reaction where they generate plasma cells and more memory B cells. It is unclear whether the memory B cells undergo further affinity maturation within these secondary GCs. In vitro activation of memory B cells can be achieved through stimulation with various activators, such as pokeweed mitogen or anti-CD40 monoclonal antibodies, however, a study found a combination of R-848 and recombinant human IL-2 to be the most efficient activator.
B cell types
Plasmablast A short-lived, proliferating antibody-secreting cell arising from B cell differentiation. Plasmablasts are generated early in an infection and their antibodies tend to have a weaker affinity towards their target antigen compared to plasma cell. Plasmablasts can result from T cell-independent activation of B cells or the extrafollicular response from T cell-dependent activation of B cells.
Plasma cell A long-lived, non-proliferating antibody-secreting cell arising from B cell differentiation. There is evidence that B cells first differentiate into a plasmablast-like cell, then differentiate into a plasma cell. Plasma cells are generated later in an infection and, compared to plasmablasts, have antibodies with a higher affinity towards their target antigen due to affinity maturation in the germinal center (GC) and produce more antibodies. Plasma cells typically result from the germinal center reaction from T cell-dependent activation of B cells, though they can also result from T cell-independent activation of B cells.
Lymphoplasmacytoid cell A cell with a mixture of B lymphocyte and plasma cell morphological features that is thought to be closely related to or a subtype of plasma cells. This cell type is found in pre-malignant and malignant plasma cell dyscrasias that are associated with the secretion of IgM monoclonal proteins; these dyscrasias include IgM monoclonal gammopathy of undetermined significance and Waldenström's macroglobulinemia.
Memory B cell Dormant B cell arising from B cell differentiation. Their function is to circulate through the body and initiate a stronger, more rapid antibody response (known as the anamnestic secondary antibody response) if they detect the antigen that had activated their parent B cell (memory B cells and their parent B cells share the same BCR, thus they detect the same antigen). Memory B cells can be generated from T cell-dependent activation through both the extrafollicular response and the germinal center reaction as well as from T cell-independent activation of B1 cells.
B-2 cell FO B cells and MZ B cells.
Follicular (FO) B cell (also known as a B-2 cell) Most common type of B cell and, when not circulating through the blood, is found mainly in the lymphoid follicles of secondary lymphoid organs (SLOs). They are responsible for generating the majority of high-affinity antibodies during an infection.
Marginal-zone (MZ) B cell Found mainly in the marginal zone of the spleen and serves as a first line of defense against blood-borne pathogens, as the marginal zone receives large amounts of blood from the general circulation. They can undergo both T cell-independent and T cell-dependent activation, but preferentially undergo T cell-independent activation.
B-1 cell Arises from a developmental pathway different from FO B cells and MZ B cells. In mice, they predominantly populate the peritoneal cavity and pleural cavity, generate natural antibodies (antibodies produced without infection), defend against mucosal pathogens, and primarily exhibit T cell-independent activation. A true homologue of mouse B-1 cells has not been discovered in humans, though various cell populations similar to B-1 cells have been described.
Regulatory B (Breg) cell An immunosuppressive B cell type that stops the expansion of pathogenic, pro-inflammatory lymphocytes through the secretion of IL-10, IL-35, and TGF-β. Also, it promotes the generation of regulatory T (Treg) cells by directly interacting with T cells to skew their differentiation towards Tregs. No common Breg cell identity has been described and many Breg cell subsets sharing regulatory functions have been found in both mice and humans. It is currently unknown if Breg cell subsets are developmentally linked and how exactly differentiation into a Breg cell occurs. There is evidence showing that nearly all B cell types can differentiate into a Breg cell through mechanisms involving inflammatory signals and BCR recognition.
B cell-related pathology
Autoimmune disease can result from abnormal B cell recognition of self-antigens followed by the production of autoantibodies. Autoimmune diseases where disease activity is correlated with B cell activity include scleroderma, multiple sclerosis, systemic lupus erythematosus, type 1 diabetes, post-infectious IBS, and rheumatoid arthritis.
Malignant transformation of B cells and their precursors can cause a host of cancers, including chronic lymphocytic leukemia (CLL), acute lymphoblastic leukemia (ALL), hairy cell leukemia, follicular lymphoma, non-Hodgkin's lymphoma, Hodgkin's lymphoma, and plasma cell malignancies such as multiple myeloma, Waldenström's macroglobulinemia, and certain forms of amyloidosis.
Abnormal B cells may be relatively large and some diseases include this in their names, such as diffuse large B-cell lymphomas (DLBCLs) and intravascular large B-cell lymphoma.
Patients with B cell alymphocytosis are predisposed to infections.
Epigenetics
A study that investigated the methylome of B cells along their differentiation cycle, using whole-genome bisulfite sequencing (WGBS), showed that there is a hypomethylation from the earliest stages to the most differentiated stages. The largest methylation difference is between the stages of germinal center B cells and memory B cells. Furthermore, this study showed that there is a similarity between B cell tumors and long-lived B cells in their DNA methylation signatures.
| Biology and health sciences | Immune system | Biology |
211955 | https://en.wikipedia.org/wiki/Natural%20killer%20cell | Natural killer cell | Natural killer cells, also known as NK cells, are a type of cytotoxic lymphocyte critical to the innate immune system. They are a kind of large granular lymphocytes (LGL), and belong to the rapidly expanding family of known innate lymphoid cells (ILC) and represent 5–20% of all circulating lymphocytes in humans. The role of NK cells is analogous to that of cytotoxic T cells in the vertebrate adaptive immune response. NK cells provide rapid responses to virus-infected cells, stressed cells, tumor cells, and other intracellular pathogens based on signals from several activating and inhibitory receptors. Most immune cells detect the antigen presented on major histocompatibility complex I (MHC-I) on infected cell surfaces, but NK cells can recognize and kill stressed cells in the absence of antibodies and MHC, allowing for a much faster immune reaction. They were named "natural killers" because of the notion that they do not require activation to kill cells that are missing "self" markers of MHC class I. This role is especially important because harmful cells that are missing MHC I markers cannot be detected and destroyed by other immune cells, such as T lymphocyte cells.
NK cells can be identified by the presence of CD56 and the absence of CD3 (CD56+, CD3−). NK cells differentiate from CD127+ common innate lymphoid progenitor, which is downstream of the common lymphoid progenitor from which B and T lymphocytes are also derived. NK cells are known to differentiate and mature in the bone marrow, lymph nodes, spleen, tonsils, and thymus, where they then enter into the circulation. NK cells differ from natural killer T cells (NKTs) phenotypically, by origin and by respective effector functions; often, NKT cell activity promotes NK cell activity by secreting interferon gamma. In contrast to NKT cells, NK cells do not express T-cell antigen receptors (TCR) or pan T marker CD3 or surface immunoglobulins (Ig) B cell receptors, but they usually express the surface markers CD16 (FcγRIII) and CD57 in humans, NK1.1 or NK1.2 in C57BL/6 mice. The NKp46 cell surface marker constitutes, at the moment, another NK cell marker of preference being expressed in both humans, several strains of mice (including BALB/c mice) and in three common monkey species.
Outside of innate immunity, both activating and inhibitory NK cell receptors play important functional roles in self tolerance and the sustaining of NK cell activity. NK cells also play a role in the adaptive immune response: numerous experiments have demonstrated their ability to readily adjust to the immediate environment and formulate antigen-specific immunological memory, fundamental for responding to secondary infections with the same antigen. The role of NK cells in both the innate and adaptive immune responses is becoming increasingly important in research using NK cell activity as a potential cancer therapy and HIV therapy.
Early history
In early experiments on cell-mediated cytotoxicity against tumor target cells, both in cancer patients and animal models, investigators consistently observed what was termed a "natural" reactivity; that is, a certain population of cells seemed to be able to destroy tumor cells without having been previously sensitized to them. The first published study to assert that untreated lymphoid cells were able to confer a natural immunity to tumors was performed by Dr. Henry Smith at the University of Leeds School of Medicine in 1966, leading to the conclusion that the "phenomenon appear[ed] to be an expression of defense mechanisms to tumor growth present in normal mice." Other researchers had also made similar observations, but as these discoveries were inconsistent with the established model at the time, many initially considered these observations to be artifacts.
By 1973, 'natural killing' activity was established across a wide variety of species, and the existence of a separate lineage of cells possessing this ability was postulated. The discovery that a unique type of lymphocyte was responsible for "natural" or spontaneous cytotoxicity was made in the early 1970s by doctoral student Rolf Kiessling and postdoctoral fellow Hugh Pross, in the mouse, and by Hugh Pross and doctoral student Mikael Jondal in the human. The mouse and human work was carried out under the supervision of professors Eva Klein and Hans Wigzell, respectively, of the Karolinska Institute, Stockholm. Kiessling's research involved the well-characterized ability of T lymphocytes to attack tumor cells which they had been previously immunized against. Pross and Jondal were studying cell-mediated cytotoxicity in normal human blood and the effect of the removal of various receptor-bearing cells on this cytotoxicity. Later that same year, Ronald Herberman published similar data with respect to the unique nature of the mouse effector cell.
The human data were confirmed, for the most part, by West et al. using similar techniques and the same erythroleukemic target cell line, K562. K562 is highly sensitive to lysis by human NK cells and, over the decades, the K562 51chromium-release assay has become the most commonly used assay to detect human NK functional activity. Its almost universal use has meant that experimental data can be compared easily by different laboratories around the world.
Using discontinuous density centrifugation, and later monoclonal antibodies, natural killing ability was mapped to the subset of large, granular lymphocytes known today as NK cells. The demonstration that density gradient-isolated large granular lymphocytes were responsible for human NK activity, made by Timonen and Saksela in 1980, was the first time that NK cells had been visualized microscopically, and was a major breakthrough in the field.
Types
NK cells can be classified as CD56bright or CD56dim. CD56bright NK cells are similar to T helper cells in exerting their influence by releasing cytokines. CD56bright NK cells constitute the majority of NK cells, being found in bone marrow, secondary lymphoid tissue, liver, and skin. CD56bright NK cells are characterized by their preferential killing of highly proliferative cells, and thus might have an immunoregulatory role. CD56dim NK cells are primarily found in the peripheral blood, and are characterized by their cell killing ability. CD56dim NK cells are always CD16 positive (CD16 is the key mediator of antibody-dependent cellular cytotoxicity, or ADCC). CD56bright can transition into CD56dim by acquiring CD16.
NK cells can eliminate virus-infected cells via CD16-mediated ADCC. All coronavirus disease 2019 (COVID-19) patients show depleted CD56bright NK cells, but CD56dim is only depleted in patients with severe COVID-19.
Receptors
NK cell receptors can also be differentiated based on function. Natural cytotoxicity receptors directly induce apoptosis (cell death) after binding to Fas ligand that directly indicate infection of a cell. The MHC-independent receptors (described above) use an alternate pathway to induce apoptosis in infected cells. Natural killer cell activation is determined by the balance of inhibitory and activating receptor stimulation. For example, if the inhibitory receptor signaling is more prominent, then NK cell activity will be inhibited; similarly, if the activating signal is dominant, then NK cell activation will result.
NK cell receptor types (with inhibitory, as well as some activating members) are differentiated by structure, with a few examples to follow:
Activating receptors
Ly49 (homodimers), relatively ancient, C-type lectin family receptors, are of multigenic presence in mice, while humans have only one pseudogenic Ly49, the receptor for classical (polymorphic) MHC I molecules.
NCR (natural cytotoxicity receptors), type 1 transmembrane proteins of the immunoglobulin superfamily, upon stimulation mediate NK killing and release of IFNγ. They bind viral ligands such as hemagglutinins and hemagglutinin neuraminidases, some bacterial ligands and cellular ligands related to tumour growth such as PCNA.
CD16 (FcγIIIA) plays a role in antibody-dependent cell-mediated cytotoxicity; in particular, they bind immunoglobulin G.
TLR – Toll-like receptors are receptors that belong in the group of pattern recognition receptors (PRR) which are typical for the cells of innate immunity but are expressed also on NK cells. They recognize PAMPs (pathogen-associated molecular patterns) and DAMPs (damage-associated molecular patterns) as their ligands. These receptors are crucial for the induction of the immune response. TLR induction amplifies the immune response by promoting the production of inflammatory cytokines and chemokines and ultimately leads to the activation of NK cell effector functions. So NK cells directly reacts to the presence of pathogens in its surroundings. Apart from TLR-10 NK cells express all of the human TLR although in various levels. NK cells express high levels of TLR-1, moderate levels of TLR-2, TLR-3, TLR-5 and TLR-6, low levels of TLR-4, TLR-8 and TLR-9 and very low levels of TLR-7. TLR receptors are constitutionally expressed independently of their state of activation and they cooperate with cytokines and chemokines on the activation of the natural killer cells. These receptors are expressed extracellularly on the cell surface or endosomally inside the endosomes. Apart from TLR-3 and TLR-4, all TLR signal through adaptor protein MyD88 which ultimately leads mainly to the activation of NF-κB. TLR-3 signals through the adaptor protein TRIF and TLR-4 can switch between signaling through MyD88 and TRIF respectively. Induction of different TLR leads to distinct activation of NK cell functions.
Inhibitory receptors
Killer-cell immunoglobulin-like receptors (KIRs) belong to a multigene family of more recently evolved Ig-like extracellular domain receptors; they are present in nonhuman primates, and are the main receptors for both classical MHC I (HLA-A, HLA-B, HLA-C) and nonclassical Mamu-G (HLA-G) in primates. Some KIRs are specific for certain HLA subtypes. Most KIRs are inhibitory and dominant. Regular cells express MHC class 1, so are recognised by KIR receptors and NK cell killing is inhibited.
CD94/NKG2 (heterodimers), a C-type lectin family receptor, is conserved in both rodents and primates and identifies nonclassical (also nonpolymorphic) MHC I molecules such as HLA-E. Expression of HLA-E at the cell surface is dependent on the presence of nonamer peptide epitope derived from the signal sequence of classical MHC class I molecules, which is generated by the sequential action of signal peptide peptidase and the proteasome. Though indirect, this is a way to survey the levels of classical (polymorphic) HLA molecules.
ILT or LIR (immunoglobulin-like receptor) – are recently discovered members of the Ig receptor family.
Ly49 (homodimers) have both activating and inhibitory isoforms. They are highly polymorphic on the population level; though they are structurally unrelated to KIRs, they are the functional homologues of KIRs in mice, including the expression pattern. Ly49s are receptor for classical (polymorphic) MHC I molecules.
Function
Cytolytic granule mediated cell apoptosis
NK cells are cytotoxic; small granules in their cytoplasm contain proteins such as perforin and proteases known as granzymes. Upon release in close proximity to a cell slated for killing, perforin forms pores in the cell membrane of the target cell, creating an aqueous channel through which the granzymes and associated molecules can enter, inducing either apoptosis or osmotic cell lysis. The distinction between apoptosis and cell lysis is important in immunology: lysing a virus-infected cell could potentially release the virions, whereas apoptosis leads to destruction of the virus inside. α-defensins, antimicrobial molecules, are also secreted by NK cells, and directly kill bacteria by disrupting their cell walls in a manner analogous to that of neutrophils.
Antibody-dependent cell-mediated cytotoxicity (ADCC)
Infected cells are routinely opsonized with antibodies for detection by immune cells. Antibodies that bind to antigens can be recognised by FcγRIII (CD16) receptors expressed on NK cells, resulting in NK activation, release of cytolytic granules and consequent cell apoptosis. This is a major killing mechanism of some monoclonal antibodies like rituximab (Rituxan), ofatumumab (Azzera), and others. The contribution of antibody-dependent cell-mediated cytotoxicity to tumor cell killing can be measured with a specific test that uses NK-92, an immortal line of NK-like cells licensed to NantKwest, Inc.: the response of NK-92 cells that have been transfected with a high-affinity Fc receptor are compared to that of the "wild type" NK-92 which does not express the Fc receptor.
Cytokine-induced NK and Cytotoxic T lymphocyte (CTL) activation
Cytokines play a crucial role in NK cell activation. As these are stress molecules released by cells upon viral infection, they serve to signal to the NK cell the presence of viral pathogens in the affected area. Cytokines involved in NK activation include IL-12, IL-15, IL-18, IL-2, and CCL5. NK cells are activated in response to interferons or macrophage-derived cytokines. They serve to contain viral infections while the adaptive immune response generates antigen-specific cytotoxic T cells that can clear the infection. NK cells work to control viral infections by secreting IFNγ and TNFα. IFNγ activates macrophages for phagocytosis and lysis, and TNFα acts to promote direct NK tumor cell killing. Patients deficient in NK cells prove to be highly susceptible to early phases of herpes virus infection. [Citation needed]
Missing 'self' hypothesis
For NK cells to defend the body against viruses and other pathogens, they require mechanisms that enable the determination of whether a cell is infected or not. The exact mechanisms remain the subject of current investigation, but recognition of an "altered self" state is thought to be involved. To control their cytotoxic activity, NK cells possess two types of surface receptors: activating receptors and inhibitory receptors, including killer-cell immunoglobulin-like receptors. Most of these receptors are not unique to NK cells and can be present in some T cell subsets, as well.
The inhibitory receptors recognize MHC class I alleles, which could explain why NK cells preferentially kill cells that possess low levels of MHC class I molecules. This mode of NK cell target interaction is known as "missing-self recognition", a term coined by Klas Kärre and co-workers in the late 90s. MHC class I molecules are the main mechanism by which cells display viral or tumor antigens to cytotoxic T cells. A common evolutionary adaptation to this is seen in both intracellular microbes and tumors: the chronic down-regulation of MHC I molecules, which makes affected cells invisible to T cells, allowing them to evade T cell-mediated immunity. NK cells apparently evolved as an evolutionary response to this adaptation (the loss of the MHC eliminates CD4/CD8 action, so another immune cell evolved to fulfill the function).
Tumor cell surveillance
Natural killer cells often lack antigen-specific cell surface receptors, so are part of innate immunity, i.e. able to react immediately with no prior exposure to the pathogen. In both mice and humans, NKs can be seen to play a role in tumor immunosurveillance by directly inducing the death of tumor cells (NKs act as cytolytic effector lymphocytes), even in the absence of surface adhesion molecules and antigenic peptides. This role of NK cells is critical to immune success particularly because T cells are unable to recognize pathogens in the absence of surface antigens. Tumor cell detection results in activation of NK cells and consequent cytokine production and release.
If tumor cells do not cause inflammation, they will also be regarded as self and will not induce a T cell response. A number of cytokines are produced by NKs, including tumor necrosis factor α (TNFα), IFNγ, and interleukin (IL-10). TNFα and IL-10 act as proinflammatory and immunosuppressors, respectively. The activation of NK cells and subsequent production of cytolytic effector cells impacts macrophages, dendritic cells, and neutrophils, which subsequently enables antigen-specific T and B cell responses. Instead of acting via antigen-specific receptors, lysis of tumor cells by NK cells is mediated by alternative receptors, including NKG2D, NKp44, NKp46, NKp30, and DNAM. NKG2D is a disulfide-linked homodimer which recognizes a number of ligands, including ULBP and MICA, which are typically expressed on tumor cells. The role of dendritic cell—NK cell interface in immunobiology have been studied and defined as critical for the comprehension of the complex immune system.
NK cells, along with macrophages and several other cell types, express the Fc receptor (FcR) molecule (FC-gamma-RIII = CD16), an activating biochemical receptor that binds the Fc portion of IgG class antibodies. This allows NK cells to target cells against which there has been a humoral response and to lyse cells through antibody-dependant cytotoxicity (ADCC). This response depends on the affinity of the Fc receptor expressed on NK cells, which can have high, intermediate, and low affinity for the Fc portion of the antibody. This affinity is determined by the amino acid in position 158 of the protein, which can be phenylalanine (F allele) or valine (V allele). Individuals with high-affinity FcRgammRIII (158 V/V allele) respond better to antibody therapy. This has been shown for lymphoma patients who received the antibody Rituxan. Patients who express the 158 V/V allele had a better antitumor response. Only 15–25% of the population expresses the 158 V/V allele. To determine the ADCC contribution of monoclonal antibodies, NK-92 cells (a "pure" NK cell line) has been transfected with the gene for the high-affinity FcR.
Clearance of senescent cells
Natural killer cells (NK cells) and macrophages play a major role in clearance of senescent cells. Natural killer cells directly kill senescent cells, and produce cytokines which activate macrophages which remove senescent cells.
Natural killer cells can use NKG2D receptors to detect senescent cells, and kill those cells using perforin pore-forming cytolytic protein. CD8+ cytotoxic T-lymphocytes also use NKG2D receptors to detect senescent cells, and promote killing similar to NK cells. For example, in patients with Parkinson's disease, levels of Natural killer cells are elevated as they degrade alpha-synuclein aggregates, destroy senescent neurons, and attenuate the neuroinflammation by leukocytes in the central nervous system.
Adaptive features of NK cells—"memory-like", "adaptive" and memory NK cells
The ability to generate memory cells following a primary infection and the consequent rapid immune activation and response to succeeding infections by the same antigen is fundamental to the role that T and B cells play in the adaptive immune response. For many years, NK cells have been considered to be a part of the innate immune system. However, recently increasing evidence suggests that NK cells can display several features that are usually attributed to adaptive immune cells (e.g. T cell responses) such as dynamic expansion and contraction of subsets, increased longevity and a form of immunological memory, characterized by a more potent response upon secondary challenge with the same antigen.
In mice, the majority of research was carried out with murine cytomegalovirus (MCMV) and in models of hapten-hypersensitivity reactions. Especially, in the MCMV model, protective memory functions of MCMV-induced NK cells were discovered
and direct recognition of the MCMV-ligand m157 by the receptor Ly49 was demonstrated to be crucial for the generation of adaptive NK cell responses.
In humans, most studies have focused on the expansion of an NK cell subset carrying the activating receptor NKG2C (KLRC2). Such expansions were observed primarily in response to human cytomegalovirus (HCMV), but also in other infections including Hantavirus, Chikungunya virus, HIV, or viral hepatitis. However, whether these virus infections trigger the expansion of adaptive NKG2C+ NK cells or whether other infections result in re-activation of latent HCMV (as suggested for hepatitis
), remains a field of study. Notably, recent research suggests that adaptive NK cells can use the activating receptor NKG2C (KLRC2) to directly bind to human cytomegalovirus-derived peptide antigens and respond to peptide recognition with activation, expansion, and differentiation, a mechanism of responding to virus infections that was previously only known for T cells of the adaptive immune system.
NK cell function in pregnancy
As the majority of pregnancies involve two parents who are not tissue-matched, successful pregnancy requires the mother's immune system to be suppressed. NK cells are thought to be an important cell type in this process. These cells are known as "uterine NK cells" (uNK cells) and they differ from peripheral NK cells. They are in the CD56bright NK cell subset, potent at cytokine secretion, but with low cytotoxic ability and relatively similar to peripheral CD56bright NK cells, with a slightly different receptor profile. These uNK cells are the most abundant leukocytes present in utero in early pregnancy, representing about 70% of leukocytes here, but from where they originate remains controversial.
These NK cells have the ability to elicit cell cytotoxicity in vitro, but at a lower level than peripheral NK cells, despite containing perforin. Lack of cytotoxicity in vivo may be due to the presence of ligands for their inhibitory receptors. Trophoblast cells downregulate HLA-A and HLA-B to defend against cytotoxic T cell-mediated death. This would normally trigger NK cells by missing self recognition; however, these cells survive. The selective retention of HLA-E (which is a ligand for NK cell inhibitory receptor NKG2A) and HLA-G (which is a ligand for NK cell inhibitory receptor KIR2DL4) by the trophoblast is thought to defend it against NK cell-mediated death.
Uterine NK cells have shown no significant difference in women with recurrent miscarriage compared with controls. However, higher peripheral NK cell percentages occur in women with recurrent miscarriages than in control groups.
NK cells secrete a high level of cytokines which help mediate their function. NK cells interact with HLA-C to produce cytokines necessary for trophoblastic proliferation. Some important cytokines they secrete include TNF-α, IL-10, IFN-γ, GM-CSF and TGF-β, among others. For example, IFN-γ dilates and thins the walls of maternal spiral arteries to enhance blood flow to the implantation site.
NK cell evasion by tumor cells
By shedding decoy NKG2D soluble ligands, tumor cells may avoid immune responses. These soluble NKG2D ligands bind to NK cell NKG2D receptors, activating a false NK response and consequently creating competition for the receptor site. This method of evasion occurs in prostate cancer. In addition, prostate cancer tumors can evade CD8 cell recognition due to their ability to downregulate expression of MHC class 1 molecules. This example of immune evasion actually highlights NK cells' importance in tumor surveillance and response, as CD8 cells can consequently only act on tumor cells in response to NK-initiated cytokine production (adaptive immune response).
Excessive NK cells
Experimental treatments with NK cells have resulted in excessive cytokine production, and even septic shock. Depletion of the inflammatory cytokine interferon gamma reversed the effect.
Applications
Anticancer therapy
Tumor-infiltrating NK cells have been reported to play a critical role in promoting drug-induced cell death in human triple-negative breast cancer. Since NK cells recognize target cells when they express nonself HLA antigens (but not self), autologous (patients' own) NK cell infusions have not shown any antitumor effects. Instead, investigators are working on using allogeneic cells from peripheral blood, which requires that all T cells be removed before infusion into the patients to remove the risk of graft versus host disease, which can be fatal. This can be achieved using an immunomagnetic column (CliniMACS). In addition, because of the limited number of NK cells in blood (only 10% of lymphocytes are NK cells), their number needs to be expanded in culture. This can take a few weeks and the yield is donor-dependent.
CAR-NK cells
Chimeric antigen receptors (CARs) are genetically modified receptors targeting cell surface antigens that provide a valuable approach to enhance effector cell efficacy. CARs induce high-affinity binding of effector cells carrying these receptors to cells expressing the target antigen, thereby lowering the threshold for cellular activation and inducing effector functions.
CAR T cells are now a fairly well-known cell therapy. However, wider use is limited by several fundamental problems: The high cost of CAR T cell therapy, which is due to the need to generate specific CAR T cells for each patient; the necessity to use only autologous T cells, due to the high risk of GvHD if allogeneic T cells are used; the inability to reinfuse CAR T cells if the patient relapses or low CAR T cell survival is observed; CAR T therapy also has a high toxicity, mainly due to IFN-γ production and subsequent induction of CRS (cytokine release syndrome) and/or neurotoxicity.
The use of CAR NK cells is not limited by the need to generate patient-specific cells, and at the same time, GvHD is not caused by NK cells, thus obviating the need for autologous cells. Toxic effects of CAR T therapy, such as CSR, have not been observed with the use of CAR NK cells. Thus, NK cells are considered an interesting "off-the-shelf" product option. Compared to CAR T cells, CAR NK cells retain unchanged expression of NK cell activating receptors. Thus, NK cells recognize and kill tumor cells even if, due to a tumor-escape strategy on tumor cells, ligand expression for the CAR receptor is downregulated.
NK cells derived from umbilical cord blood have been used to generate CAR.CD19 NK cells. These cells are capable of self-producing the cytokine IL-15, thereby enhancing autocrine/paracrine expression and persistence in vivo. Administration of these modified NK cells is not associated with the development of CSR, neurotoxicity, or GvHD.
The FT596 product is the first "Off-the-Shelf", universal, and allogenic CAR NK cellular product derived from iPSCs to be authorized for use in clinical studies in the USA. It consists of an anti-CD19 CAR optimized for NK cells with a transmembrane domain for the NKG2D activation receptor, a 2B4 costimulatory domain and a CD3ζ signaling domain. Two additional key components were added: 1) a high-affinity, non-cleavable Fc receptor CD16 (hnCD16) that enables tumor targeting and enhanced antibody-dependent cell cytotoxicity without negative regulation, combined with 2) a therapeutic monoclonal antibody targeting tumor cells and an IL-15/IL-15 receptor fusion protein (IL-15RF) promoting cytokine-independent persistence.
NK-92 cells
A more efficient way to obtain high numbers of NK cells is to expand NK-92 cells, an NK cell line with all the characteristics of highly active blood Natural Killer (NK) cells but with much broader and higher cytotoxicity. NK-92 cells grow continuously in culture and can be expanded to clinical-grade numbers in bags or bioreactors. Clinical studies have shown NK-92 cells to be safe and to exhibit anti-tumor activity in patients with lung or pancreatic cancer, melanoma, and lymphoma. When NK-92 cells originate from a patient with lymphoma, they must be irradiated prior to infusion. Efforts, however, are being made to engineer the cells to eliminate the need for irradiation. The irradiated cells maintain full cytotoxicity. NK-92 are allogeneic (from a donor different from the recipient), but in clinical studies have not been shown to elicit significant host reaction.
Unmodified NK-92 cells lack CD-16, making them unable to perform antibody-dependent cellular cytotoxicity (ADCC); however, the cells have been engineered to express a high affinity Fc-receptor (CD16A, 158V) genetically linked to IL-2 that is bound to the endoplasmic reticulum (ER). These high affinity NK-92 cells can perform ADCC and have greatly expanded therapeutic utility.
NK-92 cells have also been engineered to expressed chimeric antigen receptors (CARs), in an approach similar to that used for T cells. An example of this is an NK-92 derived cell engineered with both a CD16 and an anti-PD-L1 CAR; currently in clinical development for oncology indications. A clinical grade NK-92 variant that expresses a CAR for HER2 (ErbB2) has been generated and is in a clinical study in patients with HER2 positive glioblastoma. Several other clinical grade clones have been generated expressing the CARs for PD-L1, CD19, HER-2, and EGFR. PD-L1 targeted high affinity NK cells have been given to a number of patients with solid tumors in a phase I/II study, which is underway.
NKG2D-Fc fusion protein
In a study at Boston Children's Hospital, in coordination with Dana–Farber Cancer Institute, in which immunocompromised mice had contracted lymphomas from EBV infection, an NK-activating receptor called NKG2D was fused with a stimulatory Fc portion of the EBV antibody. The NKG2D-Fc fusion proved capable of reducing tumor growth and prolonging survival of the recipients. In a transplantation model of LMP1-fueled lymphomas, the NKG2D-Fc fusion proved capable of reducing tumor growth and prolonging survival of the recipients.
In Hodgkin lymphoma, in which the malignant Hodgkin Reed-Sternberg cells are typically HLA class I deficient, immune evasion is in part mediated by skewing towards an exhausted PD-1hi NK cell phenotype, and re-activation of these NK cells appears to be one mechanism of action induced by checkpoint-blockade.
TLR ligands
Signaling through TLR can effectively activate NK cell effector functions in vitro and in vivo. TLR ligands are then potentially able to enhance NK cell effector functions during NK cell anti-tumor immunotherapy.
Trastuzumab is a monoclonal anti-HER2 antibody that is used as a treatment of the HER2+ breast cancer. NK cells are an important part of the therapeutical effect of trastzumab as NK cells recognize the antibody coated cancer cells which induces ADCC (antibody-dependent cellular cytotoxicity) reaction. TLR ligand is used in addition to trastuzumab as a means to enhance its effect. The polysaccharide krestin, which is extracted from Trametes versicolor, is a potent ligand of TLR-2 and so activates NK cells, induces the production of IFNg and enhances the ADCC caused by recognition of trastuzumab-coated cells.
Stimulation of TLR-7 induces the expression of IFN type I and other pro-inflammatory cytokines like IL-1b, IL-6 and IL-12. Mice suffering with NK cell-sensitive lymphoma RMA-S were treated with SC1 molecule. SC1 is novel small-molecule TLR-7 agonist and its repeated administration reportedly activated NK cells in TLR-7- and IFN type I- dependent manner thus reversing the NK cell anergy which ultimately lead to lysis of the tumor.
VTX-2337 is a selective TLR-8 agonist and together with monoclonal antibody cetuximab it was used as a potential therapy for the treatment of recurrent or metastatic SCCHN. Results have shown that the NK cells had become more reactive to the treatment with cetuximab antibody upon pretreatment with VTX-2337. This indicates that the stimulation of TLR-8 and subsequent activation of inflammasome enhances the CD-16 mediated ADCC reaction in patients treated with cetuximab antibody.
NK cells play a role in controlling HIV-1 infection. TLR are potent enhancers of innate antiviral immunity and potentially can reverse HIV-1 latency. Incubation of peripheral blood mononuclear cells with novel potent TLR-9 ligand MGN1703 have resulted in enhancement of NK cell effector functions, thus significantly inhibiting the spread of HIV-1 in culture of autologous CD4+ T-cells. The stimulation of TLR-9 in NK cells induced a strong antiviral innate immune response, an increase in HIV-1 transcription (indicating the reverse in latency of the virus) and it also boosted the NK cell-mediated suppression of HIV-1 infections in autologous CD4+ T cells.
New findings
Innate resistance to HIV
Recent research suggests specific KIR-MHC class I gene interactions might control innate genetic resistance to certain viral infections, including HIV and its consequent development of AIDS. Certain HLA allotypes have been found to determine the progression of HIV to AIDS; an example is the HLA-B57 and HLA-B27 alleles, which have been found to delay progression from HIV to AIDS. This is evident because patients expressing these HLA alleles are observed to have lower viral loads and a more gradual decline in CD4+ T cells numbers. Despite considerable research and data collected measuring the genetic correlation of HLA alleles and KIR allotypes, a firm conclusion has not yet been drawn as to what combination provides decreased HIV and AIDS susceptibility.
NK cells can impose immune pressure on HIV, which had previously been described only for T cells and antibodies. HIV mutates to avoid NK cell detection.
Tissue-resident NK cells
Most of our current knowledge is derived from investigations of mouse splenic and human peripheral blood NK cells. However, in recent years tissue-resident NK cell populations have been described. These tissue-resident NK cells share transcriptional similarity to tissue-resident memory T cells described previously. However, tissue-resident NK cells are not necessarily of the memory phenotype, and in fact, the majority of the tissue-resident NK cells are functionally immature. These specialized NK-cell subsets can play a role in organ homeostasis. For example, NK cells are enriched in the human liver with a specific phenotype and take part in the control of liver fibrosis. Tissue-resident NK cells have also been identified in sites like bone marrow, spleen and more recently, in lung, intestines and lymph nodes. In these sites, tissue-resident NK cells may act as reservoir for maintaining immature NK cells in humans throughout life.
Adaptive NK cells against leukemia targets
Natural killer cells are being investigated as an emerging treatment for patients with acute myeloid leukemia (AML), and cytokine-induced memory-like NK cells have shown promise with their enhanced antileukemia functionality. It has been shown that this kind of NK cell has enhanced interferon-γ production and cytotoxicity against leukemia cell lines and primary AML blasts in patients. During a phase 1 clinical trial, five out of nine patients exhibited clinical responses to the treatment, and four patients experienced a complete remission, which suggests that these NK cells have major potential as a successful translational immunotherapy approach for patients with AML in the future.
| Biology and health sciences | Immune system | Biology |
212094 | https://en.wikipedia.org/wiki/Ballistics | Ballistics | Ballistics is the field of mechanics concerned with the launching, flight behaviour and impact effects of projectiles, especially weapon munitions such as bullets, unguided bombs, rockets and the like; the science or art of designing and accelerating projectiles so as to achieve a desired performance.
A ballistic body is a free-moving body with momentum, which can be subject to forces such as those exerted by pressurized gases from a gun barrel or a propelling nozzle, normal force by rifling, and gravity and air drag during flight.
A ballistic missile is a missile that is guided only during the relatively brief initial phase of powered flight, with the trajectory subsequently governed by the laws of classical mechanics, in contrast to (for example) a cruise missile, which is aerodynamically guided in powered flight like a fixed-wing aircraft.
History and prehistory
The earliest known ballistic projectiles were stones, spears, and the throwing stick.
The oldest evidence of stone-tipped projectiles, which may or may not have been propelled by a bow (c.f. atlatl), dating to c. 280,000 years ago, were found in Ethiopia, present day East Africa.
The oldest evidence of the use of bows to shoot arrows dates to about 10,000 years ago; it is based on pinewood arrows found in the Ahrensburg valley north of Hamburg. They had shallow grooves on the base, indicating that they were shot from a bow. The oldest bow so far recovered is about 8,000 years old, found in the Holmegård swamp in Denmark.
Archery seems to have arrived in the Americas with the Arctic small tool tradition, about 4,500 years ago.
The first devices identified as guns appeared in China around 1000 AD, and by the 12th century the technology was spreading through the rest of Asia, and into Europe by the 13th century.
After millennia of empirical development, the discipline of ballistics was initially studied and developed by Italian mathematician Niccolò Tartaglia in 1531, although he continued to use segments of straight-line motion, conventions established by the Greek philosopher Aristotle and Albert of Saxony, but with the innovation that he connected the straight lines by a circular arc. Galileo established the principle of compound motion in 1638, using the principle to derive the parabolic form of the ballistic trajectory. Ballistics was put on a solid scientific and mathematical basis by Isaac Newton, with the publication of Philosophiæ Naturalis Principia Mathematica in 1687. This gave mathematical laws of motion and gravity which for the first time made it possible to successfully predict trajectories.
The word ballistics comes from the Greek , meaning "to throw".
Projectiles
A projectile is any object projected into space (empty or not) by the exertion of a force. Although any object in motion through space (for example a thrown baseball) is a projectile, the term most commonly refers to a weapon. Mathematical equations of motion are used to analyze projectile trajectory.
Examples of projectiles include balls, arrows, bullets, artillery shells, wingless rockets, etc.
Projectile launchers
Throwing
Throwing is the launching of a projectile by hand. Although some other animals can throw, humans are unusually good throwers due to their high dexterity and good timing capabilities, and it is believed that this is an evolved trait. Evidence of human throwing dates back 2 million years. The 90 mph throwing speed found in many athletes far exceeds the speed at which chimpanzees can throw things, which is about 20 mph. This ability reflects the ability of the human shoulder muscles and tendons to store elasticity until it is needed to propel an object.
Sling
A sling is a projectile weapon typically used to throw a blunt projectile such as a stone, clay or lead "sling-bullet".
A sling has a small cradle or pouch in the middle of two lengths of cord. The sling stone is placed in the pouch. The middle finger or thumb is placed through a loop on the end of one cord, and a tab at the end of the other cord is placed between the thumb and forefinger. The sling is swung in an arc, and the tab released at a precise moment. This frees the projectile to fly to the target.
Bow
A bow is a flexible piece of material which shoots aerodynamic projectiles called arrows. The arrow is perhaps the first lethal projectile ever described in discussion of ballistics. A string joins the two ends and when the string is drawn back, the ends of the stick are flexed. When the string is released, the potential energy of the flexed stick is transformed into the velocity of the arrow. Archery is the art or sport of shooting arrows from bows.
Catapult
A catapult is a device used to launch a projectile a great distance without the aid of explosive devices – particularly various types of ancient and medieval siege engines. The catapult has been used since ancient times, because it was proven to be one of the most effective mechanisms during warfare. The word "catapult" comes from the Latin , which in turn comes from the Greek (), itself from (), "against” and (), "to toss, to hurl". Catapults were invented by the ancient Greeks.
Gun
A gun is a normally tubular weapon or other device designed to discharge projectiles or other material. The projectile may be solid, liquid, gas, or energy and may be free, as with bullets and artillery shells, or captive as with Taser probes and whaling harpoons. The means of projection varies according to design but is usually effected by the action of gas pressure, either produced through the rapid combustion of a propellant or compressed and stored by mechanical means, operating on the projectile inside an open-ended tube in the fashion of a piston. The confined gas accelerates the movable projectile down the length of the tube imparting sufficient velocity to sustain the projectile's travel once the action of the gas ceases at the end of the tube or muzzle. Alternatively, acceleration via electromagnetic field generation may be employed in which case the tube may be dispensed with and a guide rail substituted.
A weapons engineer or armourer who applies the scientific principles of ballistics to design cartridges are often called a ballistician.
Rocket
A rocket is a missile, spacecraft, aircraft or other vehicle that obtains thrust from a rocket engine. Rocket engine exhaust is formed entirely from propellants carried within the rocket before use. Rocket engines work by action and reaction. Rocket engines push rockets forward simply by throwing their exhaust backwards extremely fast.
While comparatively inefficient for low speed use, rockets are relatively lightweight and powerful, capable of generating large accelerations and of attaining extremely high speeds with reasonable efficiency. Rockets are not reliant on the atmosphere and work very well in space.
Rockets for military and recreational uses date back to at least 13th century China. Significant scientific, interplanetary and industrial use did not occur until the 20th century, when rocketry was the enabling technology for the Space Age, including setting foot on the Moon. Rockets are now used for fireworks, weaponry, ejection seats, launch vehicles for artificial satellites, human spaceflight, and space exploration.
Chemical rockets are the most common type of high performance rocket and they typically create their exhaust by the combustion of rocket propellant. Chemical rockets store a large amount of energy in an easily released form, and can be very dangerous. However, careful design, testing, construction and use minimizes risks.
Subfields
Ballistics is often broken down into the following four categories:
Internal ballistics the study of the processes originally accelerating projectiles
Transition ballistics the study of projectiles as they transition to unpowered flight
External ballistics the study of the passage of the projectile (the trajectory) in flight
Terminal ballistics the study of the projectile and its effects as it ends its flight
Internal ballistics
Internal ballistics (also interior ballistics), a sub-field of ballistics, is the study of the propulsion of a projectile.
In guns, internal ballistics covers the time from the propellant's ignition until the projectile exits the gun barrel. The study of internal ballistics is important to designers and users of firearms of all types, from small-bore rifles and pistols, to high-tech artillery.
For rocket propelled projectiles, internal ballistics covers the period during which a rocket engine is providing thrust.
Transitional ballistics
Transitional ballistics, also known as intermediate ballistics, is the study of a projectile's behavior from the time it leaves the muzzle until the pressure behind the projectile is equalized, so it lies between internal ballistics and external ballistics.
External ballistics
External ballistics is the part of the science of ballistics that deals with the behaviour of a non-powered projectile in flight.
External ballistics is frequently associated with firearms, and deals with the unpowered free-flight phase of the bullet after it exits the gun barrel and before it hits the target, so it lies between transitional ballistics and terminal ballistics.
However, external ballistics is also concerned with the free-flight of rockets and other projectiles, such as balls, arrows etc.
Terminal ballistics
Terminal ballistics is the study of the behavior and effects of a projectile when it hits its target.
Terminal ballistics is relevant both for small caliber projectiles as well as for large caliber projectiles (fired from artillery). The study of extremely high velocity impacts is still very new and is as yet mostly applied to spacecraft design.
Applications
Forensic ballistics
Forensic ballistics involves analysis of bullets and bullet impacts to determine information of use to a court or other part of a legal system. Separately from ballistics information, firearm and tool mark examinations ("ballistic fingerprinting") involve analyzing firearm, ammunition, and tool mark evidence in order to establish whether a certain firearm or tool was used in the commission of a crime.
Astrodynamics
Astrodynamics is the application of ballistics and celestial mechanics to the practical problems concerning the motion of rockets and other spacecraft. The motion of these objects is usually calculated from Newton's laws of motion and Newton's law of universal gravitation. It is a core discipline within space mission design and control.
| Physical sciences | Basics_10 | null |
212101 | https://en.wikipedia.org/wiki/Diffuse%20sky%20radiation | Diffuse sky radiation | Diffuse sky radiation is solar radiation reaching the Earth's surface after having been scattered from the direct solar beam by molecules or particulates in the atmosphere. It is also called sky radiation, the determinative process for changing the colors of the sky. Approximately 23% of direct incident radiation of total sunlight is removed from the direct solar beam by scattering into the atmosphere; of this amount (of incident radiation) about two-thirds ultimately reaches the earth as photon diffused skylight radiation.
The dominant radiative scattering processes in the atmosphere are Rayleigh scattering and Mie scattering; they are elastic, meaning that a photon of light can be deviated from its path without being absorbed and without changing wavelength.
Under an overcast sky, there is no direct sunlight, and all light results from diffused skylight radiation.
Proceeding from analyses of the aftermath of the eruption of the Philippines volcano Mount Pinatubo (in June 1991) and other studies: Diffused skylight, owing to its intrinsic structure and behavior, can illuminate under-canopy leaves, permitting more efficient total whole-plant photosynthesis than would otherwise be the case; this in stark contrast to the effect of totally clear skies with direct sunlight that casts shadows onto understory leaves and thereby limits plant photosynthesis to the top canopy layer, (see below).
Color
Earth's atmosphere scatters short-wavelength light more efficiently than that of longer wavelengths. Because its wavelengths are shorter, blue light is more strongly scattered than the longer-wavelength lights, red or green. Hence, the result that when looking at the sky away from the direct incident sunlight, the human eye perceives the sky to be blue. The color perceived is similar to that presented by a monochromatic blue (at wavelength ) mixed with white light, that is, an unsaturated blue light. The explanation of blue color by Lord Rayleigh in 1871 is a famous example of applying dimensional analysis to solving problems in physics.
Scattering and absorption are major causes of the attenuation of sunlight radiation by the atmosphere. Scattering varies as a function of the ratio of particle diameters (of particulates in the atmosphere) to the wavelength of the incident radiation. When this ratio is less than about one-tenth, Rayleigh scattering occurs. (In this case, the scattering coefficient varies inversely with the fourth power of the wavelength. At larger ratios scattering varies in a more complex fashion, as described for spherical particles by the Mie theory.) The laws of geometric optics begin to apply at higher ratios.
Daily at any global venue experiencing sunrise or sunset, most of the solar beam of visible sunlight arrives nearly tangentially to Earth's surface. Here, the path of sunlight through the atmosphere is elongated such that much of the blue or green light is scattered away from the line of perceivable visible light. This phenomenon leaves the Sun's rays, and the clouds they illuminate, abundantly orange-to-red in colors, which one sees when looking at a sunset or sunrise.
For the example of the Sun at zenith, in broad daylight, the sky is blue due to Rayleigh scattering, which also involves the diatomic gases and . Near sunset and especially during twilight, absorption by ozone () significantly contributes to maintaining blue color in the evening sky.
Under an overcast sky
There is essentially no direct sunlight under an overcast sky, so all light is then diffuse sky radiation. The flux of light is not very wavelength-dependent because the cloud droplets are larger than the light's wavelength and scatter all colors approximately equally. The light passes through the translucent clouds in a manner similar to frosted glass. The intensity ranges (roughly) from of direct sunlight for relatively thin clouds down to of direct sunlight under the extreme of thickest storm clouds.
As a part of total radiation
One of the equations for total solar radiation is:
where Hb is the beam radiation irradiance, Rb is the tilt factor for beam radiation, Hd is the diffuse radiation irradiance, Rd is the tilt factor for diffuse radiation and Rr is the tilt factor for reflected radiation.
Rb is given by:
where δ is the solar declination, Φ is the latitude, β is an angle from the horizontal and h is the solar hour angle.
Rd is given by:
and Rr by:
where ρ is the reflectivity of the surface.
Agriculture and the eruption of Mt. Pinatubo
The eruption of the Philippines volcano - Mount Pinatubo in June 1991 ejected roughly of magma and "17 million metric tons"(17 teragrams) of sulfur dioxide SO2 into the air, introducing ten times as much total SO2 as the 1991 Kuwaiti fires, mostly during the explosive Plinian/Ultra-Plinian event of June 15, 1991, creating a global stratospheric SO2 haze layer which persisted for years. This resulted in the global average temperature dropping by about . Since volcanic ash falls out of the atmosphere rapidly, the negative agricultural, effects of the eruption were largely immediate and localized to a relatively small area in close proximity to the eruption, caused by the resulting thick ash cover. Globally however, despite a several-month 5% drop in overall solar irradiation, and a reduction in direct sunlight by 30%, there was no negative impact on global agriculture. Surprisingly, a 3-4 year increase in global Agricultural productivity and forestry growth was observed, excepting boreal forest regions. The means of discovery was that initially, a mysterious drop in the rate at which carbon dioxide (CO2) was filling the atmosphere was observed, which is charted in what is known as the "Keeling Curve". This led numerous scientists to assume that the reduction was due to the lowering of Earth's temperature, and with that, a, slowdown in plant and soil respiration, indicating a deleterious impact on global agriculture from the volcanic haze layer. However upon investigation, the reduction in the rate at which carbon dioxide filled the atmosphere did not match up with the hypothesis that plant respiration rates had declined. Instead the advantageous anomaly was relatively firmly linked to an unprecedented increase in the growth/net primary production, of global plant life, resulting in the increase of the carbon sink effect of global photosynthesis. The mechanism by which the increase in plant growth was possible, was that the 30% reduction of direct sunlight can also be expressed as an increase or "enhancement" in the amount of diffuse sunlight.
The diffused skylight effect
This diffused skylight, owing to its intrinsic nature, can illuminate under-canopy leaves permitting more efficient total whole-plant photosynthesis than would otherwise be the case, and also increasing evaporative cooling, from vegetated surfaces. In stark contrast, for totally clear skies and the direct sunlight that results from it, shadows are cast onto understorey leaves, limiting plant photosynthesis to the top canopy layer. This increase in global agriculture from the volcanic haze layer also naturally results as a product of other aerosols that are not emitted by volcanoes, such, "moderately thick smoke loading" pollution, as the same mechanism, the "aerosol direct radiative effect" is behind both.
| Physical sciences | Atmospheric optics | Earth science |
212141 | https://en.wikipedia.org/wiki/Power%20station | Power station | A power station, also referred to as a power plant and sometimes generating station or generating plant, is an industrial facility for the generation of electric power. Power stations are generally connected to an electrical grid.
Many power stations contain one or more generators, rotating machine that converts mechanical power into three-phase electric power. The relative motion between a magnetic field and a conductor creates an electric current.
The energy source harnessed to turn the generator varies widely. Most power stations in the world burn fossil fuels such as coal, oil, and natural gas to generate electricity. Low-carbon power sources include nuclear power, and use of renewables such as solar, wind, geothermal, and hydroelectric.
History
In early 1871 Belgian inventor Zénobe Gramme invented a generator powerful enough to produce power on a commercial scale for industry.
In 1878, a hydroelectric power station was designed and built by William, Lord Armstrong at Cragside, England. It used water from lakes on his estate to power Siemens dynamos. The electricity supplied power to lights, heating, produced hot water, ran an elevator as well as labor-saving devices and farm buildings.
In January 1882 the world's first public coal-fired power station, the Edison Electric Light Station, was built in London, a project of Thomas Edison organized by Edward Johnson. A Babcock & Wilcox boiler powered a steam engine that drove a generator. This supplied electricity to premises in the area that could be reached through the culverts of the viaduct without digging up the road, which was the monopoly of the gas companies. The customers included the City Temple and the Old Bailey. Another important customer was the Telegraph Office of the General Post Office, but this could not be reached through the culverts. Johnson arranged for the supply cable to be run overhead, via Holborn Tavern and Newgate.
In September 1882 in New York, the Pearl Street Station was established by Edison to provide electric lighting in the lower Manhattan Island area. The station ran until destroyed by fire in 1890. The station used reciprocating steam engines to turn direct-current generators. Because of the DC distribution, the service area was small, limited by voltage drop in the feeders. In 1886 George Westinghouse began building an alternating current system that used a transformer to step up voltage for long-distance transmission and then stepped it back down for indoor lighting, a more efficient and less expensive system which is similar to modern systems. The war of the currents eventually resolved in favor of AC distribution and utilization, although some DC systems persisted to the end of the 20th century. DC systems with a service radius of a mile (kilometer) or so were necessarily smaller, less efficient of fuel consumption, and more labor-intensive to operate than much larger central AC generating stations.
AC systems used a wide range of frequencies depending on the type of load; lighting load using higher frequencies, and traction systems and heavy motor load systems preferring lower frequencies. The economics of central station generation improved greatly when unified light and power systems, operating at a common frequency, were developed. The same generating plant that fed large industrial loads during the day, could feed commuter railway systems during rush hour and then serve lighting load in the evening, thus improving the system load factor and reducing the cost of electrical energy overall. Many exceptions existed, generating stations were dedicated to power or light by the choice of frequency, and rotating frequency changers and rotating converters were particularly common to feed electric railway systems from the general lighting and power network.
Throughout the first few decades of the 20th century central stations became larger, using higher steam pressures to provide greater efficiency, and relying on interconnections of multiple generating stations to improve reliability and cost. High-voltage AC transmission allowed hydroelectric power to be conveniently moved from distant waterfalls to city markets. The advent of the steam turbine in central station service, around 1906, allowed great expansion of generating capacity. Generators were no longer limited by the power transmission of belts or the relatively slow speed of reciprocating engines, and could grow to enormous sizes. For example, Sebastian Ziani de Ferranti planned what would have reciprocating steam engine ever built for a proposed new central station, but scrapped the plans when turbines became available in the necessary size. Building power systems out of central stations required combinations of engineering skill and financial acumen in equal measure. Pioneers of central station generation include George Westinghouse and Samuel Insull in the United States, Ferranti and Charles Hesterman Merz in UK, and many others.
Thermal power stations
In thermal power stations, mechanical power is produced by a heat engine that transforms thermal energy, often from combustion of a fuel, into rotational energy. Most thermal power stations produce steam, so they are sometimes called steam power stations. Not all thermal energy can be transformed into mechanical power, according to the second law of thermodynamics; therefore, there is always heat lost to the environment. If this loss is employed as useful heat, for industrial processes or district heating, the power plant is referred to as a cogeneration power plant or CHP (combined heat-and-power) plant. In countries where district heating is common, there are dedicated heat plants called heat-only boiler stations. An important class of power stations in the Middle East uses by-product heat for the desalination of water.
The efficiency of a thermal power cycle is limited by the maximum working fluid temperature produced. The efficiency is not directly a function of the fuel used. For the same steam conditions, coal-, nuclear- and gas power plants all have the same theoretical efficiency. Overall, if a system is on constantly (base load) it will be more efficient than one that is used intermittently (peak load). Steam turbines generally operate at higher efficiency when operated at full capacity.
Besides use of reject heat for process or district heating, one way to improve overall efficiency of a power plant is to combine two different thermodynamic cycles in a combined cycle plant. Most commonly, exhaust gases from a gas turbine are used to generate steam for a boiler and a steam turbine. The combination of a "top" cycle and a "bottom" cycle produces higher overall efficiency than either cycle can attain alone.
In 2018, Inter RAO UES and State Grid planned to build an 8-GW thermal power plant, which's the largest coal-fired power plant construction project in Russia.
Classification
By heat source
Fossil-fuel power stations may also use a steam turbine generator or in the case of natural gas-fired power plants may use a combustion turbine. A coal-fired power station produces heat by burning coal in a steam boiler. The steam drives a steam turbine and generator that then produces electricity. The waste products of combustion include ash, sulfur dioxide, nitrogen oxides, and carbon dioxide. Some of the gases can be removed from the waste stream to reduce pollution.
Nuclear power plants use the heat generated in a nuclear reactor's core (by the fission process) to create steam which then operates a steam turbine and generator. About 20 percent of electric generation in the US is produced by nuclear power plants.
Geothermal power plants use steam extracted from hot underground rocks. These rocks are heated by the decay of radioactive material in the Earth's core.
Biomass-fuelled power plants may be fuelled by waste from sugar cane, municipal solid waste, landfill methane, or other forms of biomass.
In integrated steel mills, blast furnace exhaust gas is a low-cost, although low-energy-density, fuel.
Waste heat from industrial processes is occasionally concentrated enough to use for power generation, usually in a steam boiler and turbine.
Solar thermal electric plants use sunlight to boil water and produce steam which turns the generator.
Hydrogen power plants can use green hydrogen from electrolysis to help balance supply and demand from Variable renewable energy sources.
By prime mover
A prime mover is a machine that converts energy of various forms into energy of motion.
Steam turbine plants use the dynamic pressure generated by expanding steam to turn the blades of a turbine. Almost all large non-hydro plants use this system. About 90 percent of all electric power produced in the world is through use of steam turbines.
Gas turbine plants use the dynamic pressure from flowing gases (air and combustion products) to directly operate the turbine. Natural-gas fuelled (and oil fueled) combustion turbine plants can start rapidly and so are used to supply "peak" energy during periods of high demand, though at higher cost than base-loaded plants. These may be comparatively small units, and sometimes completely unmanned, being remotely operated. This type was pioneered by the UK, Princetown being the world's first, commissioned in 1959.
Combined cycle plants have both a gas turbine fired by natural gas, and a steam boiler and steam turbine which use the hot exhaust gas from the gas turbine to produce electricity. This greatly increases the overall efficiency of the plant, and many new baseload power plants are combined cycle plants fired by natural gas.
Internal combustion reciprocating engines are used to provide power for isolated communities and are frequently used for small cogeneration plants. Hospitals, office buildings, industrial plants, and other critical facilities also use them to provide backup power in case of a power outage. These are usually fuelled by diesel oil, heavy oil, natural gas, and landfill gas.
Microturbines, Stirling engine and internal combustion reciprocating engines are low-cost solutions for using opportunity fuels, such as landfill gas, digester gas from water treatment plants and waste gas from oil production.
By duty
Power plants that can be dispatched (scheduled) to provide energy to a system include:
Base load power plants run nearly continually to provide that component of system load that does not vary during a day or week. Baseload plants can be highly optimized for low fuel cost, but may not start or stop quickly during changes in system load. Examples of base-load plants would include large modern coal-fired and nuclear generating stations, or hydro plants with a predictable supply of water.
Peaking power plants meet the daily peak load, which may only be for one or two hours each day. While their incremental operating cost is always higher than base load plants, they are required to ensure security of the system during load peaks. Peaking plants include simple cycle gas turbines and reciprocating internal combustion engines, which can be started up rapidly when system peaks are predicted. Hydroelectric plants may also be designed for peaking use.
Load following power plants can economically follow the variations in the daily and weekly load, at lower cost than peaking plants and with more flexibility than baseload plants.
Non-dispatchable plants include such sources as wind and solar energy; while their long-term contribution to system energy supply is predictable, on a short-term (daily or hourly) base their energy must be used as available since generation cannot be deferred. Contractual arrangements ("take or pay") with independent power producers or system interconnections to other networks may be effectively non-dispatchable.
Cooling towers
All thermal power plants produce waste heat energy as a byproduct of the useful electrical energy produced. The amount of waste heat energy equals or exceeds the amount of energy converted into useful electricity. Gas-fired power plants can achieve as much as 65% conversion efficiency, while coal and oil plants achieve around 30–49%. The waste heat produces a temperature rise in the atmosphere, which is small compared to that produced by greenhouse-gas emissions from the same power plant. Natural draft wet cooling towers at many nuclear power plants and large fossil-fuel-fired power plants use large hyperboloid chimney-like structures (as seen in the image at the right) that release the waste heat to the ambient atmosphere by the evaporation of water.
However, the mechanical induced-draft or forced-draft wet cooling towers in many large thermal power plants, nuclear power plants, fossil-fired power plants, petroleum refineries, petrochemical plants, geothermal, biomass and waste-to-energy plants use fans to provide air movement upward through down coming water and are not hyperboloid chimney-like structures. The induced or forced-draft cooling towers are typically rectangular, box-like structures filled with a material that enhances the mixing of the upflowing air and the down-flowing water.
In areas with restricted water use, a dry cooling tower or directly air-cooled radiators may be necessary, since the cost or environmental consequences of obtaining make-up water for evaporative cooling would be prohibitive. These coolers have lower efficiency and higher energy consumption to drive fans, compared to a typical wet, evaporative cooling tower.
Air-cooled condenser (ACC)
Power plants can use an air-cooled condenser, traditionally in areas with a limited or expensive water supply. Air-cooled condensers serve the same purpose as a cooling tower (heat dissipation) without using water. They consume additional auxiliary power and thus may have a higher carbon footprint compared to a traditional cooling tower.
Once-through cooling systems
Electric companies often prefer to use cooling water from the ocean or a lake, river, or cooling pond instead of a cooling tower. This single pass or once-through cooling system can save the cost of a cooling tower and may have lower energy costs for pumping cooling water through the plant's heat exchangers. However, the waste heat can cause thermal pollution as the water is discharged. Power plants using natural bodies of water for cooling are designed with mechanisms such as fish screens, to limit intake of organisms into the cooling machinery. These screens are only partially effective and as a result billions of fish and other aquatic organisms are killed by power plants each year. For example, the cooling system at the Indian Point Energy Center in New York kills over a billion fish eggs and larvae annually. A further environmental impact is that aquatic organisms which adapt to the warmer discharge water may be injured if the plant shuts down in cold weather.
Water consumption by power stations is a developing issue.
In recent years, recycled wastewater, or grey water, has been used in cooling towers. The Calpine Riverside and the Calpine Fox power stations in Wisconsin as well as the Calpine Mankato power station in Minnesota are among these facilities.
Power from renewable energy
Power stations can generate electrical energy from renewable energy sources.
Hydroelectric power station
In a hydroelectric power station water flows through turbines using hydropower to generate hydroelectricity. Power is captured from the gravitational force of water falling through penstocks to water turbines connected to generators. The amount of power available is a combination of height and water flow. A wide range of Dams may be built to raise the water level, and create a lake for storing water.
Hydropower is produced in 150 countries, with the Asia-Pacific region generating 32 percent of global hydropower in 2010. China is the largest hydroelectricity producer, with 721 terawatt-hours of production in 2010, representing around 17 percent of domestic electricity use.
Solar
Solar energy can be turned into electricity either directly in solar cells, or in a concentrating solar power plant by focusing the light to run a heat engine.
A solar photovoltaic power plant converts sunlight into direct current electricity using the photoelectric effect. Inverters change the direct current into alternating current for connection to the electrical grid. This type of plant does not use rotating machines for energy conversion.
Solar thermal power plants use either parabolic troughs or heliostats to direct sunlight onto a pipe containing a heat transfer fluid, such as oil. The heated oil is then used to boil water into steam, which turns a turbine that drives an electrical generator. The central tower type of solar thermal power plant uses hundreds or thousands of mirrors, depending on size, to direct sunlight onto a receiver on top of a tower. The heat is used to produce steam to turn turbines that drive electrical generators.
Wind
Wind turbines can be used to generate electricity in areas with strong, steady winds, sometimes offshore. Many different designs have been used in the past, but almost all modern turbines being produced today use a three-bladed, upwind design. Grid-connected wind turbines now being built are much larger than the units installed during the 1970s. They thus produce power more cheaply and reliably than earlier models. With larger turbines (on the order of one megawatt), the blades move more slowly than older, smaller, units, which makes them less visually distracting and safer for birds.
Marine
Marine energy or marine power (also sometimes referred to as ocean energy or ocean power) refers to the energy carried by ocean waves, tides, salinity, and ocean temperature differences. The movement of water in the world's oceans creates a vast store of kinetic energy, or energy in motion. This energy can be harnessed to generate electricity to power homes, transport and industries.
The term marine energy encompasses both wave power—power from surface waves, and tidal power—obtained from the kinetic energy of large bodies of moving water. Offshore wind power is not a form of marine energy, as wind power is derived from the wind, even if the wind turbines are placed over water.
The oceans have a tremendous amount of energy and are close to many if not most concentrated populations. Ocean energy has the potential of providing a substantial amount of new renewable energy around the world.
Osmosis
Salinity gradient energy is called pressure-retarded osmosis. In this method, seawater is pumped into a pressure chamber that is at a pressure lower than the difference between the pressures of saline water and fresh water. Freshwater is also pumped into the pressure chamber through a membrane, which increases both the volume and pressure of the chamber. As the pressure differences are compensated, a turbine is spun creating energy. This method is being specifically studied by the Norwegian utility Statkraft, which has calculated that up to 25 TWh/yr would be available from this process in Norway. Statkraft has built the world's first prototype osmotic power plant on the Oslo fjord which was opened on 24 November 2009. In January 2014, however, Statkraft announced not to continue this pilot.
Biomass
Biomass energy can be produced from combustion of waste green material to heat water into steam and drive a steam turbine. Bioenergy can also be processed through a range of temperatures and pressures in gasification, pyrolysis or torrefaction reactions. Depending on the desired end product, these reactions create more energy-dense products (syngas, wood pellets, biocoal) that can then be fed into an accompanying engine to produce electricity at a much lower emission rate when compared with open burning.
Storage power stations
It is possible to store energy and produce electrical power at a later time as in pumped-storage hydroelectricity, thermal energy storage, flywheel energy storage, battery storage power station and so on.
Pumped storage
The world's largest form of storage for excess electricity, pumped-storage is a reversible hydroelectric plant. They are a net consumer of energy but provide storage for any source of electricity, effectively smoothing peaks and troughs in electricity supply and demand. Pumped storage plants typically use "spare" electricity during off peak periods to pump water from a lower reservoir to an upper reservoir. Because the pumping takes place "off peak", electricity is less valuable than at peak times. This less valuable "spare" electricity comes from uncontrolled wind power and base load power plants such as coal, nuclear and geothermal, which still produce power at night even though demand is very low. During daytime peak demand, when electricity prices are high, the storage is used for peaking power, where water in the upper reservoir is allowed to flow back to a lower reservoir through a turbine and generator. Unlike coal power stations, which can take more than 12 hours to start up from cold, a hydroelectric generator can be brought into service in a few minutes, ideal to meet a peak load demand. Two substantial pumped storage schemes are in South Africa, Palmiet Pumped Storage Scheme and another in the Drakensberg, Ingula Pumped Storage Scheme.
Typical power output
The power generated by a power station is measured in multiples of the watt, typically megawatts (106 watts) or gigawatts (109 watts). Power stations vary greatly in capacity depending on the type of power plant and on historical, geographical and economic factors. The following examples offer a sense of the scale.
Many of the largest operational onshore wind farms are located in China. As of 2022, the Roscoe Wind Farm is the largest onshore wind farm in the world, producing 8000 MW of power, followed by the Zhang Jiakou (3000 MW). As of January 2022, the Hornsea Wind Farm in United Kingdom is the largest offshore wind farm in the world at 1218 MW, followed by Walney Wind Farm in United Kingdom at 1026 MW.
In 2021, the worldwide installed capacity of power plants increased by 347 GW. Solar and wind power plant capacities rose by 80% in one year. , the largest photovoltaic (PV) power plants in the world are led by Bhadla Solar Park in India, rated at 2245 MW.
Solar thermal power stations in the U.S. have the following output:
Ivanpah Solar Power Facility is the largest of the country with an output of 392 MW
Large coal-fired, nuclear, and hydroelectric power stations can generate hundreds of megawatts to multiple gigawatts. Some examples:
The Koeberg Nuclear Power Station in South Africa has a rated capacity of 1860 megawatts.
The coal-fired Ratcliffe-on-Soar Power Station in the UK has a rated capacity of 2 gigawatts.
The Aswan Dam hydro-electric plant in Egypt has a capacity of 2.1 gigawatts.
The Three Gorges Dam hydro-electric plant in China has a capacity of 22.5 gigawatts.
Gas turbine power plants can generate tens to hundreds of megawatts. Some examples:
The Indian Queens simple-cycle, or open cycle gas turbine (OCGT), peaking power station in Cornwall UK, with a single gas turbine is rated 140 megawatts.
The Medway Power Station, a combined-cycle gas turbine (CCGT) power station in Kent, UK, with two gas turbines and one steam turbine, is rated 700 megawatts.
The rated capacity of a power station is nearly the maximum electrical power that the power station can produce.
Some power plants are run at almost exactly their rated capacity all the time, as a non-load-following base load power plant, except at times of scheduled or unscheduled maintenance.
However, many power plants usually produce much less power than their rated capacity.
In some cases a power plant produces much less power than its rated capacity because it uses an intermittent energy source.
Operators try to pull maximum available power from such power plants, because their marginal cost is practically zero, but the available power varies widely—in particular, it may be zero during heavy storms at night.
In some cases operators deliberately produce less power for economic reasons.
The cost of fuel to run a load following power plant may be relatively high, and the cost of fuel to run a peaking power plant is even higher—they have relatively high marginal costs.
Operators keep power plants turned off ("operational reserve") or running at minimum fuel consumption ("spinning reserve") most of the time.
Operators feed more fuel into load following power plants only when the demand rises above what lower-cost plants (i.e., intermittent and base load plants) can produce, and then feed more fuel into peaking power plants only when the demand rises faster than the load following power plants can follow.
Output metering
Not all of the generated power of a plant is necessarily delivered into a distribution system. Power plants typically also use some of the power themselves, in which case the generation output is classified into gross generation, and net generation.
Gross generation or gross electric output is the total amount of electricity generated by a power plant over a specific period of time. It is measured at the generating terminal and is measured in kilowatt-hours (kW·h), megawatt-hours (MW·h), gigawatt-hours (GW·h) or for the largest power plants terawatt-hours (TW·h). It includes the electricity used in the plant auxiliaries and in the transformers.
Gross generation = net generation + usage within the plant (also known as in-house loads)
Net generation is the amount of electricity generated by a power plant that is transmitted and distributed for consumer use. Net generation is less than the total gross power generation as some power produced is consumed within the plant itself to power auxiliary equipment such as pumps, motors and pollution control devices. Thus
Net generation = gross generation − usage within the plant ( in-house loads)
Operations
Operating staff at a power station have several duties. Operators are responsible for the safety of the work crews that frequently do repairs on the mechanical and electrical equipment. They maintain the equipment with periodic inspections and log temperatures, pressures and other important information at regular intervals. Operators are responsible for starting and stopping the generators depending on need. They are able to synchronize and adjust the voltage output of the added generation with the running electrical system, without upsetting the system. They must know the electrical and mechanical systems to troubleshoot problems in the facility and add to the reliability of the facility. Operators must be able to respond to an emergency and know the procedures in place to deal with it.
| Technology | Electricity generation and distribution | null |
212146 | https://en.wikipedia.org/wiki/Twinkling | Twinkling | Twinkling, also called scintillation, is a generic term for variations in apparent brightness, colour, or position of a distant luminous object viewed through a medium. If the object lies outside the Earth's atmosphere, as in the case of stars and planets, the phenomenon is termed astronomical scintillation; for objects within the atmosphere, the phenomenon is termed terrestrial scintillation. As one of the three principal factors governing astronomical seeing (the others being light pollution and cloud cover), atmospheric scintillation is defined as variations in illuminance only.
In simple terms, twinkling of stars is caused by the passing of light through different layers of a turbulent atmosphere. Most scintillation effects are caused by anomalous atmospheric refraction caused by small-scale fluctuations in air density usually related to temperature gradients. Scintillation effects are always much more pronounced near the horizon than near the zenith (directly overhead), since light rays near the horizon must have longer paths through the atmosphere before reaching the observer. Atmospheric twinkling is measured quantitatively using a scintillometer.
The effects of twinkling are reduced by using a larger receiver aperture; this effect is known as aperture averaging. Many modern large telescopes also use adaptive optical systems which precisely deform the figure of a mirror in order to compensate for scintillation.
While light from stars and other astronomical objects is likely to twinkle, twinkling usually does not cause images of planets to flicker appreciably.
Stars twinkle because they are so far from Earth that they appear as point sources of light easily disturbed by Earth's atmospheric turbulence, which acts like lenses and prisms diverting the light's path. Large astronomical objects closer to Earth, like the Moon and other planets, can be resolved as objects with observable diameters. With multiple observed points of light traversing the atmosphere, their light's deviations average out and the viewer perceives less variation in light coming from them.
| Physical sciences | Basics | Astronomy |
212193 | https://en.wikipedia.org/wiki/Coagulation | Coagulation | Coagulation, also known as clotting, is the process by which blood changes from a liquid to a gel, forming a blood clot. It results in hemostasis, the cessation of blood loss from a damaged vessel, followed by repair. The process of coagulation involves activation, adhesion and aggregation of platelets, as well as deposition and maturation of fibrin.
Coagulation begins almost instantly after an injury to the endothelium that lines a blood vessel. Exposure of blood to the subendothelial space initiates two processes: changes in platelets, and the exposure of subendothelial platelet tissue factor to coagulation factor VII, which ultimately leads to cross-linked fibrin formation. Platelets immediately form a plug at the site of injury; this is called primary hemostasis. Secondary hemostasis occurs simultaneously: additional coagulation factors beyond factor VII (listed below) respond in a cascade to form fibrin strands, which strengthen the platelet plug.
Coagulation is highly conserved throughout biology. In all mammals, coagulation involves both cellular components (platelets) and proteinaceous components (coagulation or clotting factors). The pathway in humans has been the most extensively researched and is the best understood. Disorders of coagulation can result in problems with hemorrhage, bruising, or thrombosis.
List of coagulation factors
There are 13 traditional clotting factors, as named below, and other substances necessary for coagulation:
Physiology
Physiology of blood coagulation is based on hemostasis, the normal bodily process that stops bleeding. Coagulation is a part of an integrated series of haemostatic reactions, involving plasma, platelet, and vascular components.
Hemostasis consists of four main stages:
Vasoconstriction (vasospasm or vascular spasm): Here, this refers to contraction of smooth muscles in the tunica media layer of endothelium (blood vessel wall).
Activation of platelets and platelet plug formation:
Platelet activation: Platelet activators, such as platelet activating factor and thromboxane A2, activate platelets in the bloodstream, leading to attachment of platelets' membrane receptors (e.g. glycoprotein IIb/IIIa) to extracellular matrix proteins (e.g. von Willebrand factor) on cell membranes of damaged endothelial cells and exposed collagen at the site of injury.
Platelet plug formation: The adhered platelets aggregate and form a temporary plug to stop bleeding. This process is often called "primary hemostasis".
Coagulation cascade: It is a series of enzymatic reactions that lead to the formation of a stable blood clot. The endothelial cells release substances like tissue factor, which triggers the extrinsic pathway of the coagulation cascade. This is called as "secondary hemostasis".
Fibrin clot formation: Near the end of the extrinsic pathway, after thrombin completes conversion of fibrinogen into fibrin, factor XIIIa (plasma transglutaminase; activated form of fibrin-stabilizing factor) promotes fibrin cross-linking, and subsequent stabilization of fibrin, leading to the formation of a fibrin clot (final blood clot), which temporarily seals the wound to allow wound healing until its inner part is dissolved by fibrinolytic enzymes, while the clot's outer part is shed off.
After the fibrin clot is formed, clot retraction occurs and then clot resolution starts, and these two process are together called "tertiary hemostasis". Activated platelets contract their internal actin and myosin fibrils in their cytoskeleton, which leads to shrinkage of the clot volume. Plasminogen activators, such as tissue plasminogen activator (t-PA), activate plasminogen into plasmin, which promotes lysis of the fibrin clot; this restores the flow of blood in the damaged/obstructed blood vessels.
Vasoconstriction
When there is an injury to a blood vessel, the endothelial cells can release various vasoconstrictor substances, such as endothelin and thromboxane, to induce the constriction of the smooth muscles in the vessel wall. This helps reduce blood flow to the site of injury and limits bleeding.
Platelet activation and platelet plug formation
When the endothelium is damaged, the normally isolated underlying collagen is exposed to circulating platelets, which bind directly to collagen with collagen-specific glycoprotein Ia/IIa surface receptors. This adhesion is strengthened further by von Willebrand factor (vWF), which is released from the endothelium and from platelets; vWF forms additional links between the platelets' glycoprotein Ib/IX/V and A1 domain. This localization of platelets to the extracellular matrix promotes collagen interaction with platelet glycoprotein VI. Binding of collagen to glycoprotein VI triggers a signaling cascade that results in activation of platelet integrins. Activated integrins mediate tight binding of platelets to the extracellular matrix. This process adheres platelets to the site of injury.
Activated platelets release the contents of stored granules into the blood plasma. The granules include ADP, serotonin, platelet-activating factor (PAF), vWF, platelet factor 4, and thromboxane A2 (TXA2), which, in turn, activate additional platelets. The granules' contents activate a Gq-linked protein receptor cascade, resulting in increased calcium concentration in the platelets' cytosol. The calcium activates protein kinase C, which, in turn, activates phospholipase A2 (PLA2). PLA2 then modifies the integrin membrane glycoprotein IIb/IIIa, increasing its affinity to bind fibrinogen. The activated platelets change shape from spherical to stellate, and the fibrinogen cross-links with glycoprotein IIb/IIIa aid in aggregation of adjacent platelets, forming a platelet plug and thereby completing primary hemostasis).
Coagulation cascade
The coagulation cascade of secondary hemostasis has two initial pathways which lead to fibrin formation. These are the contact activation pathway (also known as the intrinsic pathway), and the tissue factor pathway (also known as the extrinsic pathway), which both lead to the same fundamental reactions that produce fibrin. It was previously thought that the two pathways of coagulation cascade were of equal importance, but it is now known that the primary pathway for the initiation of blood coagulation is the tissue factor (extrinsic) pathway. The pathways are a series of reactions, in which a zymogen (inactive enzyme precursor) of a serine protease and its glycoprotein co-factor are activated to become active components that then catalyze the next reaction in the cascade, ultimately resulting in cross-linked fibrin. Coagulation factors are generally indicated by Roman numerals, with a lowercase a appended to indicate an active form.
The coagulation factors are generally enzymes called serine proteases, which act by cleaving downstream proteins. The exceptions are tissue factor, FV, FVIII, FXIII. Tissue factor, FV and FVIII are glycoproteins, and Factor XIII is a transglutaminase. The coagulation factors circulate as inactive zymogens.
The coagulation cascade is therefore classically divided into three pathways. The tissue factor and contact activation pathways both activate the "final common pathway" of factor X, thrombin and fibrin.
Tissue factor pathway (extrinsic)
The main role of the tissue factor (TF) pathway is to generate a "thrombin burst", a process by which thrombin, the most important constituent of the coagulation cascade in terms of its feedback activation roles, is released very rapidly. FVIIa circulates in a higher amount than any other activated coagulation factor. The process includes the following steps:
Following damage to the blood vessel, FVII leaves the circulation and comes into contact with tissue factor expressed on tissue-factor-bearing cells (stromal fibroblasts and leukocytes), forming an activated complex (TF-FVIIa).
TF-FVIIa activates FIX and FX.
FVII is itself activated by thrombin, FXIa, FXII, and FXa.
The activation of FX (to form FXa) by TF-FVIIa is almost immediately inhibited by tissue factor pathway inhibitor (TFPI).
FXa and its co-factor FVa form the prothrombinase complex, which activates prothrombin to thrombin.
Thrombin then activates other components of the coagulation cascade, including FV and FVIII (which forms a complex with FIX), and activates and releases FVIII from being bound to vWF.
FVIIIa is the co-factor of FIXa, and together they form the "tenase" complex, which activates FX; and so the cycle continues. ("Tenase" is a contraction of "ten" and the suffix "-ase" used for enzymes.)
Contact activation pathway (intrinsic)
The contact activation pathway begins with formation of the primary complex on collagen by high-molecular-weight kininogen (HMWK), prekallikrein, and FXII (Hageman factor). Prekallikrein is converted to kallikrein and FXII becomes FXIIa. FXIIa converts FXI into FXIa. Factor XIa activates FIX, which with its co-factor FVIIIa form the tenase complex, which activates FX to FXa. The minor role that the contact activation pathway has in initiating blood clot formation can be illustrated by the fact that individuals with severe deficiencies of FXII, HMWK, and prekallikrein do not have a bleeding disorder. Instead, contact activation system seems to be more involved in inflammation, and innate immunity. Despite this, interference with the pathway may confer protection against thrombosis without a significant bleeding risk.
Final common pathway
The division of coagulation in two pathways is arbitrary, originating from laboratory tests in which clotting times were measured either after the clotting was initiated by glass, the intrinsic pathway; or clotting was initiated by thromboplastin (a mix of tissue factor and phospholipids), the extrinsic pathway.
Further, the final common pathway scheme implies that prothrombin is converted to thrombin only when acted upon by the intrinsic or extrinsic pathways, which is an oversimplification. In fact, thrombin is generated by activated platelets at the initiation of the platelet plug, which in turn promotes more platelet activation.
Thrombin functions not only to convert fibrinogen to fibrin, it also activates Factors VIII and V and their inhibitor protein C (in the presence of thrombomodulin). By activating Factor XIII, covalent bonds are formed that crosslink the fibrin polymers that form from activated monomers. This stabilizes the fibrin network.
The coagulation cascade is maintained in a prothrombotic state by the continued activation of FVIII and FIX to form the tenase complex until it is down-regulated by the anticoagulant pathways.
Cell-based scheme of coagulation
A newer model of coagulation mechanism explains the intricate combination of cellular and biochemical events that occur during the coagulation process in vivo. Along with the procoagulant and anticoagulant plasma proteins, normal physiologic coagulation requires the presence of two cell types for formation of coagulation complexes: cells that express tissue factor (usually extravascular) and platelets.
The coagulation process occurs in two phases. First is the initiation phase, which occurs in tissue-factor-expressing cells. This is followed by the propagation phase, which occurs on activated platelets. The initiation phase, mediated by the tissue factor exposure, proceeds via the classic extrinsic pathway and contributes to about 5% of thrombin production. The amplified production of thrombin occurs via the classic intrinsic pathway in the propagation phase; about 95% of thrombin generated will be during this second phase.
Fibrinolysis
Eventually, blood clots are reorganized and resorbed by a process termed fibrinolysis. The main enzyme responsible for this process is plasmin, which is regulated by plasmin activators and plasmin inhibitors.
Role in immune system
The coagulation system overlaps with the immune system. Coagulation can physically trap invading microbes in blood clots. Also, some products of the coagulation system can contribute to the innate immune system by their ability to increase vascular permeability and act as chemotactic agents for phagocytic cells. In addition, some of the products of the coagulation system are directly antimicrobial. For example, beta-lysine, an amino acid produced by platelets during coagulation, can cause lysis of many Gram-positive bacteria by acting as a cationic detergent. Many acute-phase proteins of inflammation are involved in the coagulation system. In addition, pathogenic bacteria may secrete agents that alter the coagulation system, e.g. coagulase and streptokinase.
Immunohemostasis is the integration of immune activation into adaptive clot formation. Immunothrombosis is the pathological result of crosstalk between immunity, inflammation, and coagulation. Mediators of this process include damage-associated molecular patterns and pathogen-associated molecular patterns, which are recognized by toll-like receptors, triggering procoagulant and proinflammatory responses such as formation of neutrophil extracellular traps.
Cofactors
Various substances are required for the proper functioning of the coagulation cascade:
Calcium and phospholipids
Calcium and phospholipids (constituents of platelet membrane) are required for the tenase and prothrombinase complexes to function. Calcium mediates the binding of the complexes via the terminal gamma-carboxy residues on Factor Xa and Factor IXa to the phospholipid surfaces expressed by platelets, as well as procoagulant microparticles or microvesicles shed from them. Calcium is also required at other points in the coagulation cascade. Calcium ions play a major role in the regulation of coagulation cascade that is paramount in the maintenance of hemostasis. Other than platelet activation, calcium ions are responsible for complete activation of several coagulation factors, including coagulation Factor XIII.
Vitamin K
Vitamin K is an essential factor to the hepatic gamma-glutamyl carboxylase that adds a carboxyl group to glutamic acid residues on factors II, VII, IX and X, as well as Protein S, Protein C and Protein Z. In adding the gamma-carboxyl group to glutamate residues on the immature clotting factors, Vitamin K is itself oxidized. Another enzyme, Vitamin K epoxide reductase (VKORC), reduces vitamin K back to its active form. Vitamin K epoxide reductase is pharmacologically important as a target of anticoagulant drugs warfarin and related coumarins such as acenocoumarol, phenprocoumon, and dicumarol. These drugs create a deficiency of reduced vitamin K by blocking VKORC, thereby inhibiting maturation of clotting factors. Vitamin K deficiency from other causes (e.g., in malabsorption) or impaired vitamin K metabolism in disease (e.g., in liver failure) lead to the formation of PIVKAs (proteins formed in vitamin K absence), which are partially or totally non-gamma carboxylated, affecting the coagulation factors' ability to bind to phospholipid.
Regulators
Several mechanisms keep platelet activation and the coagulation cascade in check. Abnormalities can lead to an increased tendency toward thrombosis:
Protein C and Protein S
Protein C is a major physiological anticoagulant. It is a vitamin K-dependent serine protease enzyme that is activated by thrombin into activated protein C (APC). Protein C is activated in a sequence that starts with Protein C and thrombin binding to a cell surface protein thrombomodulin. Thrombomodulin binds these proteins in such a way that it activates Protein C. The activated form, along with protein S and a phospholipid as cofactors, degrades FVa and FVIIIa. Quantitative or qualitative deficiency of either (protein C or protein S) may lead to thrombophilia (a tendency to develop thrombosis). Impaired action of Protein C (activated Protein C resistance), for example by having the "Leiden" variant of Factor V or high levels of FVIII, also may lead to a thrombotic tendency.
Antithrombin
Antithrombin is a serine protease inhibitor (serpin) that degrades the serine proteases: thrombin, FIXa, FXa, FXIa, and FXIIa. It is constantly active, but its adhesion to these factors is increased by the presence of heparan sulfate (a glycosaminoglycan) or the administration of heparins (different heparinoids increase affinity to FXa, thrombin, or both). Quantitative or qualitative deficiency of antithrombin (inborn or acquired, e.g., in proteinuria) leads to thrombophilia.
Tissue factor pathway inhibitor (TFPI)
Tissue factor pathway inhibitor (TFPI) limits the action of tissue factor (TF). It also inhibits excessive TF-mediated activation of FVII and FX.
Plasmin
Plasmin is generated by proteolytic cleavage of plasminogen, a plasma protein synthesized in the liver. This cleavage is catalyzed by tissue plasminogen activator (t-PA), which is synthesized and secreted by endothelium. Plasmin proteolytically cleaves fibrin into fibrin degradation products that inhibit excessive fibrin formation.
Prostacyclin
Prostacyclin (PGI2) is released by endothelium and activates platelet Gs protein-linked receptors. This, in turn, activates adenylyl cyclase, which synthesizes cAMP. cAMP inhibits platelet activation by decreasing cytosolic levels of calcium and, by doing so, inhibits the release of granules that would lead to activation of additional platelets and the coagulation cascade.
Medical assessment
Numerous medical tests are used to assess the function of the coagulation system:
Common: aPTT, PT (also used to determine INR), fibrinogen testing (often by the Clauss fibrinogen assay), platelet count, platelet function testing (often by PFA-100), thrombodynamics test.
Other: TCT, bleeding time, mixing test (whether an abnormality corrects if the patient's plasma is mixed with normal plasma), coagulation factor assays, antiphospholipid antibodies, D-dimer, genetic tests (e.g. factor V Leiden, prothrombin mutation G20210A), dilute Russell's viper venom time (dRVVT), miscellaneous platelet function tests, thromboelastography (TEG or Sonoclot), euglobulin lysis time (ELT).
The contact activation (intrinsic) pathway is initiated by activation of the contact activation system, and can be measured by the activated partial thromboplastin time (aPTT) test.
The tissue factor (extrinsic) pathway is initiated by release of tissue factor (a specific cellular lipoprotein), and can be measured by the prothrombin time (PT) test. PT results are often reported as ratio (INR value) to monitor dosing of oral anticoagulants such as warfarin.
The quantitative and qualitative screening of fibrinogen is measured by the thrombin clotting time (TCT). Measurement of the exact amount of fibrinogen present in the blood is generally done using the Clauss fibrinogen assay. Many analysers are capable of measuring a "derived fibrinogen" level from the graph of the Prothrombin time clot.
If a coagulation factor is part of the contact activation or tissue factor pathway, a deficiency of that factor will affect only one of the tests: Thus hemophilia A, a deficiency of factor VIII, which is part of the contact activation pathway, results in an abnormally prolonged aPTT test but a normal PT test. Deficiencies of common pathway factors prothrombin, fibrinogen, FX, and FV will prolong both aPTT and PT. If an abnormal PT or aPTT is present, additional testing will occur to determine which (if any) factor is present as aberrant concentrations.
Deficiencies of fibrinogen (quantitative or qualitative) will prolong PT, aPTT, thrombin time, and reptilase time.
Role in disease
Coagulation defects may cause hemorrhage or thrombosis, and occasionally both, depending on the nature of the defect.
Platelet disorders
Platelet disorders are either congenital or acquired. Examples of congenital platelet disorders are Glanzmann's thrombasthenia, Bernard–Soulier syndrome (abnormal glycoprotein Ib-IX-V complex), gray platelet syndrome (deficient alpha granules), and delta storage pool deficiency (deficient dense granules). Most are rare. They predispose to hemorrhage. Von Willebrand disease is due to deficiency or abnormal function of von Willebrand factor, and leads to a similar bleeding pattern; its milder forms are relatively common.
Decreased platelet numbers (thrombocytopenia) is due to insufficient production (e.g., myelodysplastic syndrome or other bone marrow disorders), destruction by the immune system (immune thrombocytopenic purpura), or consumption (e.g., thrombotic thrombocytopenic purpura, hemolytic-uremic syndrome, paroxysmal nocturnal hemoglobinuria, disseminated intravascular coagulation, heparin-induced thrombocytopenia). An increase in platelet count is called thrombocytosis, which may lead to formation of thromboembolisms; however, thrombocytosis may be associated with increased risk of either thrombosis or hemorrhage in patients with myeloproliferative neoplasm.
Coagulation factor disorders
The best-known coagulation factor disorders are the hemophilias. The three main forms are hemophilia A (factor VIII deficiency), hemophilia B (factor IX deficiency or "Christmas disease") and hemophilia C (factor XI deficiency, mild bleeding tendency).
Von Willebrand disease (which behaves more like a platelet disorder except in severe cases), is the most common hereditary bleeding disorder and is characterized as being inherited autosomal recessive or dominant. In this disease, there is a defect in von Willebrand factor (vWF), which mediates the binding of glycoprotein Ib (GPIb) to collagen. This binding helps mediate the activation of platelets and formation of primary hemostasis.
In acute or chronic liver failure, there is insufficient production of coagulation factors, possibly increasing risk of bleeding during surgery.
Thrombosis is the pathological development of blood clots. These clots may break free and become mobile, forming an embolus or grow to such a size that occludes the vessel in which it developed. An embolism is said to occur when the thrombus (blood clot) becomes a mobile embolus and migrates to another part of the body, interfering with blood circulation and hence impairing organ function downstream of the occlusion. This causes ischemia and often leads to ischemic necrosis of tissue. Most cases of venous thrombosis are due to acquired states (older age, surgery, cancer, immobility). Unprovoked venous thrombosis may be related to inherited thrombophilias (e.g., factor V Leiden, antithrombin deficiency, and various other genetic deficiencies or variants), particularly in younger patients with family history of thrombosis; however, thrombotic events are more likely when acquired risk factors are superimposed on the inherited state.
Pharmacology
Procoagulants
The use of adsorbent chemicals, such as zeolites, and other hemostatic agents are also used for sealing severe injuries quickly (such as in traumatic bleeding secondary to gunshot wounds). Thrombin and fibrin glue are used surgically to treat bleeding and to thrombose aneurysms. Hemostatic Powder Spray TC-325 is used to treated gastrointestinal bleeding.
Desmopressin is used to improve platelet function by activating arginine vasopressin receptor 1A.
Coagulation factor concentrates are used to treat hemophilia, to reverse the effects of anticoagulants, and to treat bleeding in people with impaired coagulation factor synthesis or increased consumption. Prothrombin complex concentrate, cryoprecipitate and fresh frozen plasma are commonly used coagulation factor products. Recombinant activated human factor VII is sometimes used in the treatment of major bleeding.
Tranexamic acid and aminocaproic acid inhibit fibrinolysis and lead to a de facto reduced bleeding rate. Before its withdrawal, aprotinin was used in some forms of major surgery to decrease bleeding risk and the need for blood products.
Anticoagulants
Anticoagulants and anti-platelet agents (together "antithrombotics") are amongst the most commonly used medications. Anti-platelet agents include aspirin, dipyridamole, ticlopidine, clopidogrel, ticagrelor and prasugrel; the parenteral glycoprotein IIb/IIIa inhibitors are used during angioplasty. Of the anticoagulants, warfarin (and related coumarins) and heparin are the most commonly used. Warfarin affects the vitamin K-dependent clotting factors (II, VII, IX, X) and protein C and protein S, whereas heparin and related compounds increase the action of antithrombin on thrombin and factor Xa. A newer class of drugs, the direct thrombin inhibitors, is under development; some members are already in clinical use (such as lepirudin, argatroban, bivalirudin and dabigatran). Also in clinical use are other small molecular compounds that interfere directly with the enzymatic action of particular coagulation factors (the directly acting oral anticoagulants: dabigatran, rivaroxaban, apixaban, and edoxaban).
History
Initial discoveries
Theories on the coagulation of blood have existed since antiquity. Physiologist Johannes Müller (1801–1858) described fibrin, the substance of a thrombus. Its soluble precursor, fibrinogen, was thus named by Rudolf Virchow (1821–1902), and isolated chemically by Prosper Sylvain Denis (1799–1863). Alexander Schmidt suggested that the conversion from fibrinogen to fibrin is the result of an enzymatic process, and labeled the hypothetical enzyme "thrombin" and its precursor "prothrombin". Arthus discovered in 1890 that calcium was essential in coagulation. Platelets were identified in 1865, and their function was elucidated by Giulio Bizzozero in 1882.
The theory that thrombin is generated by the presence of tissue factor was consolidated by Paul Morawitz in 1905. At this stage, it was known that thrombokinase/thromboplastin (factor III) is released by damaged tissues, reacting with prothrombin (II), which, together with calcium (IV), forms thrombin, which converts fibrinogen into fibrin (I).
Coagulation factors
The remainder of the biochemical factors in the process of coagulation were largely discovered in the 20th century.
A first clue as to the actual complexity of the system of coagulation was the discovery of proaccelerin (initially and later called Factor V) by (1905–1990) in 1947. He also postulated its function to be the generation of accelerin (Factor VI), which later turned out to be the activated form of V (or Va); hence, VI is not now in active use.
Factor VII (also known as serum prothrombin conversion accelerator or proconvertin, precipitated by barium sulfate) was discovered in a young female patient in 1949 and 1951 by different groups.
Factor VIII turned out to be deficient in the clinically recognized but etiologically elusive hemophilia A; it was identified in the 1950s and is alternatively called antihemophilic globulin due to its capability to correct hemophilia A.
Factor IX was discovered in 1952 in a young patient with hemophilia B named Stephen Christmas (1947–1993). His deficiency was described by Dr. Rosemary Biggs and Professor R.G. MacFarlane in Oxford, UK. The factor is, hence, called Christmas Factor. Christmas lived in Canada and campaigned for blood transfusion safety until succumbing to transfusion-related AIDS at age 46. An alternative name for the factor is plasma thromboplastin component, given by an independent group in California.
Hageman factor, now known as factor XII, was identified in 1955 in an asymptomatic patient with a prolonged bleeding time named of John Hageman. Factor X, or Stuart-Prower factor, followed, in 1956. This protein was identified in a Ms. Audrey Prower of London, who had a lifelong bleeding tendency. In 1957, an American group identified the same factor in a Mr. Rufus Stuart. Factors XI and XIII were identified in 1953 and 1961, respectively.
The view that the coagulation process is a "cascade" or "waterfall" was enunciated almost simultaneously by MacFarlane in the UK and by Davie and Ratnoff in the US, respectively.
Nomenclature
The usage of Roman numerals rather than eponyms or systematic names was agreed upon during annual conferences (starting in 1955) of hemostasis experts. In 1962, consensus was achieved on the numbering of factors I–XII. This committee evolved into the present-day International Committee on Thrombosis and Hemostasis (ICTH). Assignment of numerals ceased in 1963 after the naming of Factor XIII. The names Fletcher Factor and Fitzgerald Factor were given to further coagulation-related proteins, namely prekallikrein and high-molecular-weight kininogen, respectively.
Factor VI is unassigned, as accelerin was found to be activated Factor V.
Other species
All mammals have an extremely closely related blood coagulation process, using a combined cellular and serine protease process. It is possible for any mammalian coagulation factor to "cleave" its equivalent target in any other mammal. The only non-mammalian animal known to use serine proteases for blood coagulation is the horseshoe crab. Exemplifying the close links between coagulation and inflammation, the horseshoe crab has a primitive response to injury, carried out by cells known as amoebocytes (or hemocytes) which serve both hemostatic and immune functions.
| Biology and health sciences | Basics | Biology |
212240 | https://en.wikipedia.org/wiki/Blood%20plasma | Blood plasma | Blood plasma is a light amber-colored liquid component of blood in which blood cells are absent, but which contains proteins and other constituents of whole blood in suspension. It makes up about 55% of the body's total blood volume. It is the intravascular part of extracellular fluid (all body fluid outside cells). It is mostly water (up to 95% by volume), and contains important dissolved proteins (6–8%; e.g., serum albumins, globulins, and fibrinogen), glucose, clotting factors, electrolytes (, , , , , etc.), hormones, carbon dioxide (plasma being the main medium for excretory product transportation), and oxygen. It plays a vital role in an intravascular osmotic effect that keeps electrolyte concentration balanced and protects the body from infection and other blood-related disorders.
Blood plasma can be separated from whole blood through blood fractionation, by adding an anticoagulant to a tube filled with blood, which is spun in a centrifuge until the blood cells fall to the bottom of the tube. The blood plasma is then poured or drawn off. For point-of-care testing applications, plasma can be extracted from whole blood via filtration or via agglutination to allow for rapid testing of specific biomarkers. Blood plasma has a density of approximately . Blood serum is blood plasma without clotting factors. Plasmapheresis is a medical therapy that involves blood plasma extraction, treatment, and reintegration.
Fresh frozen plasma is on the WHO Model List of Essential Medicines, the most important medications needed in a basic health system. It is of critical importance in the treatment of many types of trauma which result in blood loss, and is therefore kept stocked universally in all medical facilities capable of treating trauma (e.g., trauma centers, hospitals, and ambulances) or that pose a risk of patient blood loss such as surgical suite facilities.
Volume
Blood plasma volume may be expanded by or drained to extravascular fluid when there are changes in Starling forces across capillary walls. For example, when blood pressure drops in circulatory shock, Starling forces drive fluid into the interstitium, causing third spacing.
Standing still for a prolonged period will cause an increase in transcapillary hydrostatic pressure. As a result, approximately 12% of blood plasma volume will cross into the extravascular compartment. This plasma shift causes an increase in hematocrit, serum total protein, blood viscosity and, as a result of increased concentration of coagulation factors, it causes orthostatic hypercoagulability.
Plasma proteins
Albumins
Serum albumins are the most common plasma proteins, and they are responsible for maintaining the osmotic pressure of the blood. Without albumins, the consistency of blood would be closer to that of water. The increased viscosity of blood prevents fluid from entering the bloodstream from outside the capillaries. Albumins are produced in the liver, assuming the absence of a hepatocellular deficiency.
Globulins
The second most common type of protein in the blood plasma are globulins. Important globulins include immunoglobins which are important for the immune system and transport hormones and other compounds around the body. There are three main types of globulins. Alpha-1 and Alpha-2 globulins are formed in the liver and play an important role in mineral transport and the inhibition of blood coagulation. An example of beta globulin found in blood plasma includes low-density lipoproteins (LDL) which are responsible for transporting fat to the cells for steroid and membrane synthesis. Gamma globulin, better known as immunoglobulins, are produced by plasma B cells, and provides the human body with a defense system against invading pathogens and other immune diseases.
Fibrinogen
Fibrinogen proteins make up most of the remaining proteins in the blood. Fibrinogens are responsible for clotting blood to help prevent blood loss.
Color
Plasma is normally yellow due to bilirubin, carotenoids, hemoglobin, and transferrin. In abnormal cases, plasma can have varying shades of orange, green, or brown. The green color can be due to ceruloplasmin or sulfhemoglobin. The latter may form due to medicines that are able to form sulfonamides once ingested. A dark brown or reddish color can appear due to hemolysis, in which methemoglobin is released from broken blood cells. Plasma is normally relatively transparent, but sometimes it can be opaque. Opaqueness is typically due to elevated content of lipids like cholesterol and triglycerides.
Plasma vs. serum in medical diagnostics
Plasma and serum are both derived from full blood, but serum is obtained by removing blood cells, fibrin clots, and other coagulation factors while plasma is obtained by only removing blood cells. Blood plasma and blood serum are often used in blood tests. Tests can be done on plasma, serum or both. In addition, some tests have to be done with whole blood, such as the determination of the amount of blood cells in blood via flow cytometry.
Benefits of plasma over serum
Plasma preparation is quick, as it is not coagulated. Serum sample preparation requires about 30 minutes of waiting time before it can be centrifuged and then analyzed. However, coagulation can be hastened down to a few minutes by adding thrombin or similar agents to the serum sample.
Compared to serum, 15–20% larger volume of plasma can be obtained from a blood sample of certain size. Serum lacks some proteins that partake in coagulation and increase the sample volume.
Serum preparation can cause measurement errors by increasing or decreasing the concentration of the analyte that is meant to be measured. For example, during coagulation, blood cells consume blood glucose and platelets increase the sample content of compounds like potassium, phosphates and aspartate transaminase by secreting them. Glucose or these other compounds may be the analytes.
Benefits of serum over plasma
Plasma preparation requires the addition of anticoagulants, which can cause expected and unexpected measurement errors. For example, anticoagulant salts can add extra cations like NH4+, Li+, Na+ and K+ to the sample, or impurities like lead and aluminum. Chelator anticoagulants like EDTA and citrate salts work by binding calcium (see carboxyglutamic acid), but they may also bind other ions. Even if such ions are not the analytes, chelators can interfere with enzyme activity measurements. For example, EDTA binds zinc ions, which alkaline phosphatases need as cofactors. Thus, phosphatase activity cannot be measured if EDTA is used.
An unknown volume of anticoagulants can be added to a plasma sample by accident, which may ruin the sample as the analyte concentration is changed by an unknown amount.
No anticoagulants are added to serum samples, which decreases the preparation cost of the samples relative to plasma samples.
Plasma samples can form tiny clots if the added anticoagulant is not properly mixed with the sample. Non-uniform samples can cause measurement errors.
History
Plasma was already well known when described by William Harvey in de Motu Cordis in 1628, but knowledge of it probably dates as far back as Vesalius (1514–1564). The discovery of fibrinogen by William Henson, , made it easier to study plasma, as ordinarily, upon coming in contact with a foreign surface – something other than the vascular endothelium – clotting factors become activated and clotting proceeds rapidly, trapping RBCs etc. in the plasma and preventing separation of plasma from the blood. Adding citrate and other anticoagulants is a relatively recent advance. Upon the formation of a clot, the remaining clear fluid (if any) is blood serum, which is essentially plasma without the clotting factors
The use of blood plasma as a substitute for whole blood and for transfusion purposes was proposed in March 1918, in the correspondence columns of the British Medical Journal, by Gordon R. Ward. "Dried plasmas" in powder or strips of material format were developed and first used in World War II. Prior to the United States' involvement in the war, liquid plasma and whole blood were used.
The origin of plasmapheresis
Dr. José Antonio Grifols Lucas, a scientist from Vilanova i la Geltrú, Spain, founded Laboratorios Grifols in 1940. Dr. Grifols pioneered a first-of-its-kind technique called plasmapheresis, where a donor's red blood cells would be returned to the donor's body almost immediately after the separation of the blood plasma. This technique is still in practice today, almost 80 years later. In 1945, Dr. Grifols opened the world's first plasma donation center.
Blood for Britain
The "Blood for Britain" program during the early 1940s was quite successful (and popular in the United States) based on Charles Drew's contribution. A large project began in August 1940 to collect blood in New York City hospitals for the export of plasma to Britain. Drew was appointed medical supervisor of the "Plasma for Britain" project. His notable contribution at this time was to transform the test tube methods of many blood researchers into the first successful mass production techniques.
Nevertheless, the decision was made to develop a dried plasma package for the armed forces as it would reduce breakage and make the transportation, packaging, and storage much simpler. The resulting dried plasma package came in two tin cans containing 400 cc bottles. One bottle contained enough distilled water to reconstitute the dried plasma contained within the other bottle. In about three minutes, the plasma would be ready to use and could stay fresh for around four hours. The Blood for Britain program operated successfully for five months, with total collections of almost 15,000 people donating blood, and with over 5,500 vials of blood plasma.
Following the Supplying Blood Plasma to England project, Drew was named director of the Red Cross blood bank and assistant director of the National Research Council, in charge of blood collection for the United States Army and Navy. Drew argued against the armed forces directive that blood/plasma was to be separated by the race of the donor. Drew insisted that there was no racial difference in human blood and that the policy would lead to needless deaths as soldiers and sailors were required to wait for "same race" blood.
By the end of the war the American Red Cross had provided enough blood for over six million plasma packages. Most of the surplus plasma was returned to the United States for civilian use. Serum albumin replaced dried plasma for combat use during the Korean War.
Plasma donation
Plasma as a blood product prepared from blood donations is used in blood transfusions, typically as fresh frozen plasma (FFP) or Plasma Frozen within 24 hours after phlebotomy (PF24). When donating whole blood or packed red blood cell (PRBC) transfusions, O- is the most desirable and is considered a "universal donor," since it has neither A nor B antigens and can be safely transfused to most recipients. Type AB+ is the "universal recipient" type for PRBC donations. However, for plasma the situation is somewhat reversed. Blood donation centers will sometimes collect only plasma from AB donors through apheresis, as their plasma does not contain the antibodies that may cross react with recipient antigens. As such, AB is often considered the "universal donor" for plasma. Special programs exist just to cater to the male AB plasma donor, because of concerns about transfusion related acute lung injury (TRALI) and female donors who may have higher leukocyte antibodies. However, some studies show an increased risk of TRALI despite increased leukocyte antibodies in women who have been pregnant.
United Kingdom
Following fears of variant Creutzfeldt-Jakob disease (vCJD) being spread through the blood supply, the British government began to phase out blood plasma from U.K. donors and by the end of 1999 had imported all blood products made with plasma from the United States. In 2002, the British government purchased Life Resources Incorporated, an American blood supply company, to import plasma. The company became Plasma Resources UK (PRUK) which owned Bio Products Laboratory. In 2013, the British government sold an 80% stake in PRUK to American hedge fund Bain Capital, in a deal estimated to be worth £200 million. The sale was met with criticism in the UK. In 2009, the U.K. stopped importing plasma from the United States, as it was no longer a viable option due to regulatory and jurisdictional challenges.
At present (2024), blood donated in the United Kingdom is used by UK Blood Services for the manufacture of plasma blood components (Fresh Frozen Plasma (FFP) and cryoprecipitate). However, plasma from UK donors is still not used for the commercial manufacture of fractionated plasma medicines.
Synthetic blood plasma
Simulated body fluid (SBF) is a solution having a similar ion concentration to that of human blood plasma. SBF is normally used for the surface modification of metallic implants, and more recently in gene delivery application.
| Biology and health sciences | Circulatory system | Biology |
212250 | https://en.wikipedia.org/wiki/Homotopy | Homotopy | In topology, two continuous functions from one topological space to another are called homotopic (from "same, similar" and "place") if one can be "continuously deformed" into the other, such a deformation being called a homotopy (, ; , ) between the two functions. A notable use of homotopy is the definition of homotopy groups and cohomotopy groups, important invariants in algebraic topology.
In practice, there are technical difficulties in using homotopies with certain spaces. Algebraic topologists work with compactly generated spaces, CW complexes, or spectra.
Formal definition
Formally, a homotopy between two continuous functions f and g from a
topological space X to a topological space Y is defined to be a continuous function from the product of the space X with the unit interval [0, 1] to Y such that and for all .
If we think of the second parameter of H as time then H describes a continuous deformation of f into g: at time 0 we have the function f and at time 1 we have the function g. We can also think of the second parameter as a "slider control" that allows us to smoothly transition from f to g as the slider moves from 0 to 1, and vice versa.
An alternative notation is to say that a homotopy between two continuous functions is a family of continuous functions for such that and , and the map is continuous from to . The two versions coincide by setting . It is not sufficient to require each map to be continuous.
The animation that is looped above right provides an example of a homotopy between two embeddings, f and g, of the torus into . X is the torus, Y is , f is some continuous function from the torus to R3 that takes the torus to the embedded surface-of-a-doughnut shape with which the animation starts; g is some continuous function that takes the torus to the embedded surface-of-a-coffee-mug shape. The animation shows the image of ht(X) as a function of the parameter t, where t varies with time from 0 to 1 over each cycle of the animation loop. It pauses, then shows the image as t varies back from 1 to 0, pauses, and repeats this cycle.
Properties
Continuous functions f and g are said to be homotopic if and only if there is a homotopy H taking f to g as described above. Being homotopic is an equivalence relation on the set of all continuous functions from X to Y.
This homotopy relation is compatible with function composition in the following sense: if are homotopic, and are homotopic, then their compositions and are also homotopic.
Examples
If are given by and , then the map given by is a homotopy between them.
More generally, if is a convex subset of Euclidean space and are paths with the same endpoints, then there is a linear homotopy (or straight-line homotopy) given by
Let be the identity function on the unit n-disk; i.e. the set . Let be the constant function which sends every point to the origin. Then the following is a homotopy between them:
Homotopy equivalence
Given two topological spaces X and Y, a homotopy equivalence between X and Y is a pair of continuous maps and , such that is homotopic to the identity map idX and is homotopic to idY. If such a pair exists, then X and Y are said to be homotopy equivalent, or of the same homotopy type. Intuitively, two spaces X and Y are homotopy equivalent if they can be transformed into one another by bending, shrinking and expanding operations. Spaces that are homotopy-equivalent to a point are called contractible.
Homotopy equivalence vs. homeomorphism
A homeomorphism is a special case of a homotopy equivalence, in which is equal to the identity map idX (not only homotopic to it), and is equal to idY. Therefore, if X and Y are homeomorphic then they are homotopy-equivalent, but the opposite is not true. Some examples:
A solid disk is homotopy-equivalent to a single point, since you can deform the disk along radial lines continuously to a single point. However, they are not homeomorphic, since there is no bijection between them (since one is an infinite set, while the other is finite).
The Möbius strip and an untwisted (closed) strip are homotopy equivalent, since you can deform both strips continuously to a circle. But they are not homeomorphic.
Examples
The first example of a homotopy equivalence is with a point, denoted . The part that needs to be checked is the existence of a homotopy between and , the projection of onto the origin. This can be described as .
There is a homotopy equivalence between (the 1-sphere) and .
More generally, .
Any fiber bundle with fibers homotopy equivalent to a point has homotopy equivalent total and base spaces. This generalizes the previous two examples since is a fiber bundle with fiber .
Every vector bundle is a fiber bundle with a fiber homotopy equivalent to a point.
for any , by writing as the total space of the fiber bundle , then applying the homotopy equivalences above.
If a subcomplex of a CW complex is contractible, then the quotient space is homotopy equivalent to .
A deformation retraction is a homotopy equivalence.
Null-homotopy
A function is said to be null-homotopic if it is homotopic to a constant function. (The homotopy from to a constant function is then sometimes called a null-homotopy.) For example, a map from the unit circle to any space is null-homotopic precisely when it can be continuously extended to a map from the unit disk to that agrees with on the boundary.
It follows from these definitions that a space is contractible if and only if the identity map from to itself—which is always a homotopy equivalence—is null-homotopic.
Invariance
Homotopy equivalence is important because in algebraic topology many concepts are homotopy invariant, that is, they respect the relation of homotopy equivalence. For example, if X and Y are homotopy equivalent spaces, then:
X is path-connected if and only if Y is.
X is simply connected if and only if Y is.
The (singular) homology and cohomology groups of X and Y are isomorphic.
If X and Y are path-connected, then the fundamental groups of X and Y are isomorphic, and so are the higher homotopy groups. (Without the path-connectedness assumption, one has π1(X, x0) isomorphic to π1(Y, f(x0)) where is a homotopy equivalence and
An example of an algebraic invariant of topological spaces which is not homotopy-invariant is compactly supported homology (which is, roughly speaking, the homology of the compactification, and compactification is not homotopy-invariant).
Variants
Relative homotopy
In order to define the fundamental group, one needs the notion of homotopy relative to a subspace. These are homotopies which keep the elements of the subspace fixed. Formally: if f and g are continuous maps from X to Y and K is a subset of X, then we say that f and g are homotopic relative to K if there exists a homotopy between f and g such that for all and Also, if g is a retraction from X to K and f is the identity map, this is known as a strong deformation retract of X to K.
When K is a point, the term pointed homotopy is used.
Isotopy
When two given continuous functions f and g from the topological space X to the topological space Y are embeddings, one can ask whether they can be connected 'through embeddings'. This gives rise to the concept of isotopy, which is a homotopy, H, in the notation used before, such that for each fixed t, H(x, t) gives an embedding.
A related, but different, concept is that of ambient isotopy.
Requiring that two embeddings be isotopic is a stronger requirement than that they be homotopic. For example, the map from the interval [−1, 1] into the real numbers defined by f(x) = −x is not isotopic to the identity g(x) = x. Any homotopy from f to the identity would have to exchange the endpoints, which would mean that they would have to 'pass through' each other. Moreover, f has changed the orientation of the interval and g has not, which is impossible under an isotopy. However, the maps are homotopic; one homotopy from f to the identity is H: [−1, 1] × [0, 1] → [−1, 1] given by H(x, y) = 2yx − x.
Two homeomorphisms (which are special cases of embeddings) of the unit ball which agree on the boundary can be shown to be isotopic using Alexander's trick. For this reason, the map of the unit disc in defined by f(x, y) = (−x, −y) is isotopic to a 180-degree rotation around the origin, and so the identity map and f are isotopic because they can be connected by rotations.
In geometric topology—for example in knot theory—the idea of isotopy is used to construct equivalence relations. For example, when should two knots be considered the same? We take two knots, K1 and K2, in three-dimensional space. A knot is an embedding of a one-dimensional space, the "loop of string" (or the circle), into this space, and this embedding gives a homeomorphism between the circle and its image in the embedding space. The intuitive idea behind the notion of knot equivalence is that one can deform one embedding to another through a path of embeddings: a continuous function starting at t = 0 giving the K1 embedding, ending at t = 1 giving the K2 embedding, with all intermediate values corresponding to embeddings. This corresponds to the definition of isotopy. An ambient isotopy, studied in this context, is an isotopy of the larger space, considered in light of its action on the embedded submanifold. Knots K1 and K2 are considered equivalent when there is an ambient isotopy which moves K1 to K2. This is the appropriate definition in the topological category.
Similar language is used for the equivalent concept in contexts where one has a stronger notion of equivalence. For example, a path between two smooth embeddings is a smooth isotopy.
Timelike homotopy
On a Lorentzian manifold, certain curves are distinguished as timelike (representing something that only goes forwards, not backwards, in time, in every local frame). A timelike homotopy between two timelike curves is a homotopy such that the curve remains timelike during the continuous transformation from one curve to another. No closed timelike curve (CTC) on a Lorentzian manifold is timelike homotopic to a point (that is, null timelike homotopic); such a manifold is therefore said to be multiply connected by timelike curves. A manifold such as the 3-sphere can be simply connected (by any type of curve), and yet be timelike multiply connected.
Properties
Lifting and extension properties
If we have a homotopy and a cover and we are given a map such that (h0 is called a lift of h0), then we can lift all H to a map such that The homotopy lifting property is used to characterize fibrations.
Another useful property involving homotopy is the homotopy extension property,
which characterizes the extension of a homotopy between two functions from a subset of some set to the set itself. It is useful when dealing with cofibrations.
Groups
Since the relation of two functions being homotopic relative to a subspace is an equivalence relation, we can look at the equivalence classes of maps between a fixed X and Y. If we fix , the unit interval [0, 1] crossed with itself n times, and we take its boundary as a subspace, then the equivalence classes form a group, denoted , where is in the image of the subspace .
We can define the action of one equivalence class on another, and so we get a group. These groups are called the homotopy groups. In the case , it is also called the fundamental group.
Homotopy category
The idea of homotopy can be turned into a formal category of category theory. The homotopy category is the category whose objects are topological spaces, and whose morphisms are homotopy equivalence classes of continuous maps. Two topological spaces X and Y are isomorphic in this category if and only if they are homotopy-equivalent. Then a functor on the category of topological spaces is homotopy invariant if it can be expressed as a functor on the homotopy category.
For example, homology groups are a functorial homotopy invariant: this means that if f and g from X to Y are homotopic, then the group homomorphisms induced by f and g on the level of homology groups are the same: Hn(f) = Hn(g) : Hn(X) → Hn(Y) for all n. Likewise, if X and Y are in addition path connected, and the homotopy between f and g is pointed, then the group homomorphisms induced by f and g on the level of homotopy groups are also the same: πn(f) = πn(g) : πn(X) → πn(Y).
Applications
Based on the concept of the homotopy, computation methods for algebraic and differential equations have been developed. The methods for algebraic equations include the homotopy continuation method and the continuation method (see numerical continuation). The methods for differential equations include the homotopy analysis method.
Homotopy theory can be used as a foundation for homology theory: one can represent a cohomology functor on a space X by mappings of X into an appropriate fixed space, up to homotopy equivalence. For example, for any abelian group G, and any based CW-complex X, the set of based homotopy classes of based maps from X to the Eilenberg–MacLane space is in natural bijection with the n-th singular cohomology group of the space X. One says that the omega-spectrum of Eilenberg-MacLane spaces are representing spaces for singular cohomology with coefficients in G.
| Mathematics | Geometry | null |
212253 | https://en.wikipedia.org/wiki/Electric%20power%20distribution | Electric power distribution | Electric power distribution is the final stage in the delivery of electricity. Electricity is carried from the transmission system to individual consumers. Distribution substations connect to the transmission system and lower the transmission voltage to medium voltage ranging between and with the use of transformers. Primary distribution lines carry this medium voltage power to distribution transformers located near the customer's premises. Distribution transformers again lower the voltage to the utilization voltage used by lighting, industrial equipment and household appliances. Often several customers are supplied from one transformer through secondary distribution lines. Commercial and residential customers are connected to the secondary distribution lines through service drops. Customers demanding a much larger amount of power may be connected directly to the primary distribution level or the subtransmission level.
The transition from transmission to distribution happens in a power substation, which has the following functions:
Circuit breakers and switches enable the substation to be disconnected from the transmission grid or for distribution lines to be disconnected.
Transformers step down transmission voltages, or more, down to primary distribution voltages. These are medium voltage circuits, usually .
From the transformer, power goes to the busbar that can split the distribution power off in multiple directions. The bus distributes power to distribution lines, which fan out to customers.
Urban distribution is mainly underground, sometimes in common utility ducts. Rural distribution is mostly above ground with utility poles, and suburban distribution is a mix.
Closer to the customer, a distribution transformer steps the primary distribution power down to a low-voltage secondary circuit, usually 120/240 V in the US for residential customers. The power comes to the customer via a service drop and an electricity meter. The final circuit in an urban system may be less than but may be over for a rural customer.
History
Electric power distribution become necessary only in the 1880s, when electricity started being generated at power stations. Until then, electricity was usually generated where it was used. The first power-distribution systems installed in European and US cities were used to supply lighting: arc lighting running on very-high-voltage (around 3,000 V) alternating current (AC) or direct current (DC), and incandescent lighting running on low-voltage (100 V) direct current. Both were supplanting gas lighting systems, with arc lighting taking over large-area and street lighting, and incandescent lighting replacing gas lights for business and residential users.
The high voltages used in arc lighting allowed a single generating station to supply a string of lights up to long. And each doubling of voltage would allow a given cable to transmit the same amount of power four times the distance than at the lower voltage (with the same power loss). By contrast, direct-current indoor incandescent lighting systems, such as Edison's first power station, installed in 1882, had difficulty supplying customers more than a mile away because they used a low voltage (110 V) from generation to end use. The low voltage translated to higher current and required thick copper cables for transmission. In practice, Edison's DC generating plants needed to be within about of the farthest customer to avoid even thicker and more expensive conductors.
Introduction of the transformer
The problem of transmitting electricity over longer distances became a recognized engineering roadblock to electric power distribution, with many less-than-satisfactory solutions tested by lighting companies. But the mid-1880s saw a breakthrough with the development of functional transformers that allowed AC power to be "stepped up" to a much higher voltage for transmission, then dropped down to a lower voltage near the end user. Compared to direct current, AC had much cheaper transmission costs and greater economies of scale — with large AC generating plants capable of supplying whole cities and regions, which led to the use of AC spreading rapidly.
In the US the competition between direct current and alternating current took a personal turn in the late 1880s in the form of a "war of currents" when Thomas Edison started attacking George Westinghouse and his development of the first US AC transformer systems, highlighting the deaths caused by high-voltage AC systems over the years and claiming any AC system was inherently dangerous. Edison's propaganda campaign was short-lived, with his company switching over to AC in 1892.
AC became the dominant form of transmission of power with innovations in Europe and the US in electric motor designs, and the development of engineered universal systems allowing the large number of legacy systems to be connected to large AC grids.
In the first half of the 20th century, in many places the electric power industry was vertically integrated, meaning that one company did generation, transmission, distribution, metering and billing. Starting in the 1970s and 1980s, nations began the process of deregulation and privatization, leading to electricity markets. The distribution system would remain regulated, but generation, retail, and sometimes transmission systems were transformed into competitive markets.
Generation and transmission
Electric power begins at a generating station, where the potential difference can be as high as 33,000 volts. AC is usually used. Users of large amounts of DC power such as some railway electrification systems, telephone exchanges and industrial processes such as aluminium smelting use rectifiers to derive DC from the public AC supply, or may have their own generation systems. High-voltage DC can be advantageous for isolating alternating-current systems or controlling the quantity of electricity transmitted. For example, Hydro-Québec has a direct-current line which goes from the James Bay region to Boston.
From the generating station it goes to the generating station's switchyard where a step-up transformer increases the voltage to a level suitable for transmission, from 44 kV to 765 kV. Once in the transmission system, electricity from each generating station is combined with electricity produced elsewhere. For alternating-current generators, all generating units connected to a common network must be synchronized, operating at the same frequency within a small tolerance. Alternatively, disparate sources can be combined to serve a common load if some external power converter, such as a rotating machine or a direct current converter system is interposed. Electricity is consumed as soon as it is produced. It is transmitted at a very high speed, close to the speed of light.
Primary distribution
Primary distribution voltages range from 4 kV to 35 kV phase-to-phase (2.4 kV to 20 kV phase-to-neutral) Only large consumers are fed directly from distribution voltages; most utility customers are connected to a transformer, which reduces the distribution voltage to the low voltage "utilization voltage", "supply voltage" or "mains voltage" used by lighting and interior wiring systems.
Network configurations
Distribution networks are divided into two types, radial or network. A radial system is arranged like a tree where each customer has one source of supply. A network system has multiple sources of supply operating in parallel. Spot networks are used for concentrated loads. Radial systems are commonly used in rural or suburban areas.
Radial systems usually include emergency connections where the system can be reconfigured in case of problems, such as a fault or planned maintenance. This can be done by opening and closing switches to isolate a certain section from the grid.
Long feeders experience voltage drop (power factor distortion) requiring capacitors or voltage regulators to be installed.
Reconfiguration, by exchanging the functional links between the elements of the system, represents one of the most important measures which can improve the operational performance of a distribution system. The problem of optimization through the reconfiguration of a power distribution system, in terms of its definition, is a historical single objective problem with constraints. Since 1975, when Merlin and Back introduced the idea of distribution system reconfiguration for active power loss reduction, until nowadays, a lot of researchers have proposed diverse methods and algorithms to solve the reconfiguration problem as a single objective problem. Some authors have proposed Pareto optimality based approaches (including active power losses and reliability indices as objectives). For this purpose, different artificial intelligence based methods have been used: microgenetic, branch exchange, particle swarm optimization and non-dominated sorting genetic algorithm.
Rural services
Rural electrification systems tend to use higher distribution voltages because of the longer distances covered by distribution lines (see Rural Electrification Administration). 7.2, 12.47, 25, and 34.5 kV distribution is common in the United States; 11 kV and 33 kV are common in the UK, Australia and New Zealand; 11 kV and 22 kV are common in South Africa; 10, 20 and 35 kV are common in China. Other voltages are occasionally used.
Rural services normally try to minimize the number of poles and wires. It uses higher voltages (than urban distribution), which in turn permits use of galvanized steel wire. The strong steel wire allows for less expensive wide pole spacing. In rural areas a pole-mount transformer may serve only one customer. In New Zealand, Australia, Saskatchewan, Canada, and South Africa, Single-wire earth return systems (SWER) are used to electrify remote rural areas.
Three phase service provides power for large agricultural facilities, petroleum pumping facilities, water plants, or other customers that have large loads (three-phase equipment).
In North America, overhead distribution systems may be three phase, four wire, with a neutral conductor. Rural distribution system may have long runs of one phase conductor and a neutral.
In other countries or in extreme rural areas the neutral wire is connected to the ground to use that as a return (single-wire earth return).
Secondary distribution
Electricity is delivered at a frequency of either 50 or 60 Hz, depending on the region. It is delivered to domestic customers as single-phase electric power. In some countries as in Europe a three phase supply may be made available for larger properties. Seen with an oscilloscope, the domestic power supply in North America would look like a sine wave, oscillating between −170 volts and 170 volts, giving an effective voltage of 120 volts RMS. Three-phase electric power is more efficient in terms of power delivered per cable used, and is more suited to running large electric motors. Some large European appliances may be powered by three-phase power, such as electric stoves and clothes dryers.
A ground connection is normally provided for the customer's system as well as for the equipment owned by the utility. The purpose of connecting the customer's system to ground is to limit the voltage that may develop if high voltage conductors fall down onto lower-voltage conductors which are usually mounted lower to the ground, or if a failure occurs within a distribution transformer. Earthing systems can be TT, TN-S, TN-C-S or TN-C.
Regional variations
220–240 volt systems
Most of the world uses 50 Hz 220 or 230 V single phase, or 400 V three-phase for residential and light industrial services. In this system, the primary distribution network supplies a few substations per area, and the 230 V / 400 V power from each substation is directly distributed to end users over a region of normally less than 1 km radius. Three live (hot) wires and the neutral are connected to the building for a three phase service. Single-phase distribution, with one live wire and the neutral is used domestically where total loads are light. In Europe, electricity is normally distributed for industry and domestic use by the three-phase, four wire system. This gives a phase-to-phase voltage of wye service and a single-phase voltage of between any one phase and neutral. In the UK a typical urban or suburban low-voltage substation would normally be rated between 150 kVA and 1 MVA and supply a whole neighbourhood of a few hundred houses. Transformers are typically sized on an average load of 1 to 2 kW per household, and the service fuses and cable is sized to allow any one property to draw a peak load of perhaps ten times this.
For industrial customers, 3-phase is also available, or may be generated locally.
Large industrial customers have their own transformer(s) with an input from 11 kV to 220 kV.
100–120 volt systems
Most of the Americas use 60 Hz AC, the 120/240 volt split-phase system domestically and three phase for larger installations. North American transformers usually power homes at 240 volts, similar to Europe's 230 volts. It is the split-phase that allows use of 120 volts in the home.
In the electricity sector in Japan, the standard voltage is 100 V, with both 50 and 60 Hz AC frequencies being used. Parts of the country use 50 Hz, while other parts use 60 Hz. This is a relic from the 1890s. Some local providers in Tokyo imported 50 Hz German equipment, while the local power providers in Osaka brought in 60 Hz generators from the United States. The grids grew until eventually the entire country was wired. Today the frequency is 50 Hz in Eastern Japan (including Tokyo, Yokohama, Tohoku, and Hokkaido) and 60 Hz in Western Japan (including Nagoya, Osaka, Kyoto, Hiroshima, Shikoku, and Kyushu).
Most household appliances are made to work on either frequency. The problem of incompatibility came into the public eye when the 2011 Tōhoku earthquake and tsunami knocked out about a third of the east's capacity, and power in the west could not be fully shared with the east since the country does not have a common frequency.
There are four high-voltage direct current (HVDC) converter stations that move power across Japan's AC frequency border. Shin Shinano is a back-to-back HVDC facility in Japan which forms one of four frequency changer stations that link Japan's western and eastern power grids. The other three are at Higashi-Shimizu, Minami-Fukumitsu and Sakuma Dam. Together they can move up to 1.2 GW of power east or west.
240 volt systems and 120 volt outlets
Most modern North American homes are wired to receive 240 volts from the transformer, and through the use of split-phase electrical power, can have both 120 volt receptacles and 240 volt receptacles. The 120 volts is typically used for lighting and most wall outlets. The 240 volt circuits are typically used for appliances requiring high watt heat output such as ovens and heaters. They may also be used to supply an electric car charger.
Modern distribution systems
Traditionally, the distribution systems would only operate as simple distribution lines where the electricity from the transmission networks would be shared among the customers. Today's distribution systems are heavily integrated with renewable energy generations at the distribution level of the power systems by the means of distributed generation resources, such as solar energy and wind energy. As a result, distribution systems are becoming more independent from the transmission networks day-by-day. Balancing the supply-demand relationship at these modern distribution networks (sometimes referred to as microgrids) is extremely challenging, and it requires the use of various technological and operational means to operate. Such tools include battery storage power station, data analytics, optimization tools, etc.
| Technology | Electricity generation and distribution | null |
212330 | https://en.wikipedia.org/wiki/Lye | Lye | Lye is a hydroxide, either sodium hydroxide or potassium hydroxide. The word lye most accurately refers to sodium hydroxide (NaOH), but historically has been conflated to include other alkali materials, most notably potassium hydroxide (KOH). In order to distinguish between the two, sodium hydroxide may be referred to as soda lye while potassium hydroxide may be referred to as potash lye.
Traditionally, it was obtained by using rainwater to leach wood ashes (which are highly soluble in water and strongly alkaline) of their potassium hydroxide (KOH). A caustic basic solution is produced, called lye water. Then, the lye water would either be used as such, as for curing olives before brining them, or be evaporated of water to produce crystalline lye.
Today, lye is commercially manufactured using a membrane cell chloralkali process. It is supplied in various forms such as flakes, pellets, microbeads, coarse powder or a solution. Lye has traditionally been used as a major ingredient in soapmaking.
Etymology
The English word has cognates in all Germanic languages, and originally designated a bath or hot spring.
Uses
Food
Lyes are used to cure many types of food, including the traditional Nordic lutefisk, olives (making them less bitter), canned mandarin oranges, hominy, lye rolls, century eggs, pretzels, candied pumpkins, and bagels. They are also used as a tenderizer in the crust of baked Cantonese moon cakes, in "zongzi" (glutinous rice dumplings wrapped in bamboo leaves), in chewy southern Chinese noodles popular in Hong Kong and southern China, and in Japanese ramen noodles. Lye provides the crisp glaze on hard pretzels. It's used in kutsinta, a type of rice cake from the Philippines together with pitsi-pitsî. In Assam, north east India, extensive use is made of a type of lye called khar in Assamese and karwi in Boro which is obtained by filtering the ashes of various banana stems, roots and skin in their cooking and also for curing, as medicine and as a substitute for soap. Lye made out of wood ashes is also used in the nixtamalization process of hominy corn by the tribes of the Eastern Woodlands in North America.
In the United States, food-grade lye must meet the requirements outlined in the Food Chemicals Codex (FCC), as prescribed by the U.S. Food and Drug Administration (FDA). Lower grades of lye that are unsuitable for use in food preparation are commonly used as drain cleaners and oven cleaners.
Soap
Both sodium hydroxide and potassium hydroxide are used in making soap. Potassium hydroxide soaps are softer and more easily dissolved in water than sodium hydroxide soaps. Sodium hydroxide and potassium hydroxide are not interchangeable in either the proportions required or the properties produced in making soaps.
"Hot process" soap making also uses lye as the main ingredient. Lye is added to water, cooled for a few minutes and then added to oils and butters. The mixture is then cooked over a period of time (1–2 hours), typically in a slow cooker, and then placed into a mold.
Household
Lyes are also valued for their cleaning effects. Sodium hydroxide is commonly the major constituent in commercial and industrial oven cleaners and clogged drain openers, due to its grease-dissolving abilities. Lyes decompose greases via alkaline ester hydrolysis, yielding water-soluble residues that are easily removed by rinsing.
Tissue digestion
Sodium or potassium hydroxide can be used to digest tissues of animal carcasses. Often referred to as alkaline hydrolysis, the process involves placing the animal carcass into a sealed chamber, adding a mixture of lye and water and the application of heat to accelerate the process. After several hours the chamber will contain a liquid with coffee-like appearance, and the only solids that remain are very fragile bone hulls of mostly calcium phosphate, which can be mechanically crushed to a fine powder with very little force. Sodium hydroxide is frequently used in the process of decomposing roadkill dumped in landfills by animal disposal contractors. Due to its low cost and easy availability, it has also been used to dispose of corpses by criminals. Italian serial killer Leonarda Cianciulli used this chemical to turn dead bodies into soap. In Mexico, a man who worked for drug cartels admitted to having disposed of more than 300 bodies with it.
Fungus identification
A 3–10% solution of potassium hydroxide (KOH) gives a color change in some species of mushrooms:
In Agaricus, some species such as A. xanthodermus turn yellow with KOH, many have no reaction, and A. subrutilescens turns green.
Distinctive change occurs for some species of Cortinarius and boletes
Safety
First aid
When a person has been exposed to lye, sources recommend immediate removal of contaminated clothing/materials, gently brushing/wiping excess off of skin, and then flushing the area of exposure with running water for 15–60 minutes as well as contacting emergency services.
Protection
Personal protective equipment including safety glasses, chemical-resistant gloves, and adequate ventilation are required for the safe handling of lye. When in proximity to lye that is dissolving in an open container of water, the use of a vapor-resistant face mask is recommended. Adding lye too quickly can cause a runaway thermal reaction which can result in the mixture boiling or erupting.
Storage
Lye in its solid state is deliquescent and has a strong affinity for moisture in the air. As a result, lye will dissolve when exposed to open air, absorbing large amounts of atmospheric moisture. Accordingly, lye is stored in air-tight (and correspondingly moisture tight) containers. Glass is not a good material to be used for storage as severe alkalis are mildly corrosive to it. Similar to the case of other corrosives, the containers should be labeled to indicate the potential danger of the contents and stored away from children, pets, heat, and moisture.
Hazardous reactions
The majority of safety concerns with lye are also common with most corrosives, such as their potentially destructive effects on living tissues; examples are the skin, flesh, and the cornea. Solutions containing lyes can cause chemical burns, permanent injuries, scarring and blindness, immediately upon contact. Lyes may be harmful or even fatal if swallowed; ingestion can cause esophageal stricture. Moreover, the solvation of dry solid lye is highly exothermic and the resulting heat may cause additional burns or ignite flammables.
The reaction between sodium hydroxide and some metals is also hazardous. Aluminium, magnesium, zinc, tin, chromium, brass and bronze all react with lye to produce hydrogen gas. Since hydrogen is flammable, mixing a large quantity of lye with aluminium could result in an explosion. Both the potassium and sodium forms are able to dissolve copper.
| Physical sciences | Specific bases | Chemistry |
212374 | https://en.wikipedia.org/wiki/Nebular%20hypothesis | Nebular hypothesis | The nebular hypothesis is the most widely accepted model in the field of cosmogony to explain the formation and evolution of the Solar System (as well as other planetary systems). It suggests the Solar System is formed from gas and dust orbiting the Sun which clumped up together to form the planets. The theory was developed by Immanuel Kant and published in his Universal Natural History and Theory of the Heavens (1755) and then modified in 1796 by Pierre Laplace. Originally applied to the Solar System, the process of planetary system formation is now thought to be at work throughout the universe. The widely accepted modern variant of the nebular theory is the solar nebular disk model (SNDM) or solar nebular model. It offered explanations for a variety of properties of the Solar System, including the nearly circular and coplanar orbits of the planets, and their motion in the same direction as the Sun's rotation. Some elements of the original nebular theory are echoed in modern theories of planetary formation, but most elements have been superseded.
According to the nebular theory, stars form in massive and dense clouds of molecular hydrogen—giant molecular clouds (GMC). These clouds are gravitationally unstable, and matter coalesces within them to smaller denser clumps, which then rotate, collapse, and form stars. Star formation is a complex process, which always produces a gaseous protoplanetary disk (proplyd) around the young star. This may give birth to planets in certain circumstances, which are not well known. Thus the formation of planetary systems is thought to be a natural result of star formation. A Sun-like star usually takes approximately 1 million years to form, with the protoplanetary disk evolving into a planetary system over the next 10–100 million years.
The protoplanetary disk is an accretion disk that feeds the central star. Initially very hot, the disk later cools in what is known as the T Tauri star stage; here, formation of small dust grains made of rocks and ice is possible. The grains eventually may coagulate into kilometer-sized planetesimals. If the disk is massive enough, the runaway accretions begin, resulting in the rapid—100,000 to 300,000 years—formation of Moon- to Mars-sized planetary embryos. Near the star, the planetary embryos go through a stage of violent mergers, producing a few terrestrial planets. The last stage takes approximately 100 million to a billion years.
The formation of giant planets is a more complicated process. It is thought to occur beyond the frost line, where planetary embryos mainly are made of various types of ice. As a result, they are several times more massive than in the inner part of the protoplanetary disk. What follows after the embryo formation is not completely clear. Some embryos appear to continue to grow and eventually reach 5–10 Earth masses—the threshold value, which is necessary to begin accretion of the hydrogen–helium gas from the disk. The accumulation of gas by the core is initially a slow process, which continues for several million years, but after the forming protoplanet reaches about 30 Earth masses () it accelerates and proceeds in a runaway manner. Jupiter- and Saturn-like planets are thought to accumulate the bulk of their mass during only 10,000 years. The accretion stops when the gas is exhausted. The formed planets can migrate over long distances during or after their formation. Ice giants such as Uranus and Neptune are thought to be failed cores, which formed too late when the disk had almost disappeared.
History
There is evidence that Emanuel Swedenborg first proposed parts of the nebular theory in 1734. Immanuel Kant, familiar with Swedenborg's work, developed the theory further in 1755, publishing his own Universal Natural History and Theory of the Heavens, wherein he argued that gaseous clouds (nebulae) slowly rotate, gradually collapse and flatten due to gravity, eventually forming stars and planets.
Pierre-Simon Laplace independently developed and proposed a similar model in 1796 in his Exposition du systeme du monde. He envisioned that the Sun originally had an extended hot atmosphere throughout the volume of the Solar System. His theory featured a contracting and cooling protosolar cloud—the protosolar nebula. As this cooled and contracted, it flattened and spun more rapidly, throwing off (or shedding) a series of gaseous rings of material; and according to him, the planets condensed from this material. His model was similar to Kant's, except more detailed and on a smaller scale. While the Laplacian nebular model dominated in the 19th century, it encountered a number of difficulties. The main problem involved angular momentum distribution between the Sun and planets. The planets have 99% of the angular momentum, and this fact could not be explained by the nebular model. As a result, astronomers largely abandoned this theory of planet formation at the beginning of the 20th century.
According to some, a major critique came during the 19th century from James Clerk Maxwell (1831–1879), who in some sources is claimed to have maintained that different rotation between the inner and outer parts of a ring could not allow condensation of material. However, both the critique and the attribution to Maxwell have been deemed to be incorrect upon further investigation, with the original error being made by George Gamow in some popular publications and propagated continually ever since. Astronomer Sir David Brewster also rejected Laplace, writing in 1876 that "those who believe in the Nebular Theory consider it as certain that our Earth derived its solid matter and its atmosphere from a ring thrown from the Solar atmosphere, which afterwards contracted into a solid terraqueous sphere, from which the Moon was thrown off by the same process". He argued that under such view, "the Moon must necessarily have carried off water and air from the watery and aerial parts of the Earth and must have an atmosphere". Brewster claimed that Sir Isaac Newton's religious beliefs had previously considered nebular ideas as tending to atheism, and quoted him as saying that "the growth of new systems out of old ones, without the mediation of a Divine power, seemed to him apparently absurd".
The perceived deficiencies of the Laplacian model stimulated scientists to find a replacement for it. During the 20th century many theories addressed the issue, including the planetesimal theory of Thomas Chamberlin and Forest Moulton (1901), the tidal model of James Jeans (1917), the accretion model of Otto Schmidt (1944), the protoplanet theory of William McCrea (1960) and finally the capture theory of Michael Woolfson. In 1978 Andrew Prentice resurrected the initial Laplacian ideas about planet formation and developed the modern Laplacian theory. None of these attempts proved completely successful, and many of the proposed theories were descriptive.
The birth of the modern widely accepted theory of planetary formation—the solar nebular disk model (SNDM)—can be traced to the Soviet astronomer Victor Safronov. His 1969 book Evolution of the protoplanetary cloud and formation of the Earth and the planets, which was translated to English in 1972, had a long-lasting effect on the way scientists think about the formation of the planets. In this book almost all major problems of the planetary formation process were formulated and some of them solved. Safronov's ideas were further developed in the works of George Wetherill, who discovered runaway accretion. While originally applied only to the Solar System, the SNDM was subsequently thought by theorists to be at work throughout the Universe; as of astronomers have discovered extrasolar planets in our galaxy.
Solar nebular model: achievements and problems
Achievements
The star formation process naturally results in the appearance of accretion disks around young stellar objects. At the age of about 1 million years, 100% of stars may have such disks. This conclusion is supported by the discovery of the gaseous and dusty disks around protostars and T Tauri stars as well as by theoretical considerations. Observations of these disks show that the dust grains inside them grow in size on short (thousand-year) time scales, producing 1 centimeter sized particles.
The accretion process, by which 1 km planetesimals grow into 1,000 km sized bodies, is well understood now. This process develops inside any disk where the number density of planetesimals is sufficiently high, and proceeds in a runaway manner. Growth later slows and continues as oligarchic accretion. The end result is formation of planetary embryos of varying sizes, which depend on the distance from the star. Various simulations have demonstrated that the merger of embryos in the inner part of the protoplanetary disk leads to the formation of a few Earth-sized bodies. Thus the origin of terrestrial planets is now considered to be an almost solved problem.
Current issues
The physics of accretion disks encounters some problems. The most important one is how the material, which is accreted by the protostar, loses its angular momentum. One possible explanation suggested by Hannes Alfvén was that angular momentum was shed by the solar wind during its T Tauri star phase. The momentum is transported to the outer parts of the disk by viscous stresses. Viscosity is generated by macroscopic turbulence, but the precise mechanism that produces this turbulence is not well understood. Another possible process for shedding angular momentum is magnetic braking, where the spin of the star is transferred into the surrounding disk via that star's magnetic field. The main processes responsible for the disappearance of the gas in disks are viscous diffusion and photo-evaporation.
The formation of planetesimals is the biggest unsolved problem in the nebular disk model. How 1 cm sized particles coalesce into 1 km planetesimals is a mystery. This mechanism appears to be the key to the question as to why some stars have planets, while others have nothing around them, not even dust belts.
The formation timescale of giant planets is also an important problem. Old theories were unable to explain how their cores could form fast enough to accumulate significant amounts of gas from the quickly disappearing protoplanetary disk. The mean lifetime of the disks, which is less than ten million (107) years, appeared to be shorter than the time necessary for the core formation. Much progress has been done to solve this problem and current models of giant planet formation are now capable of forming Jupiter (or more massive planets) in about 4 million years or less, well within the average lifetime of gaseous disks.
Another potential problem of giant planet formation is their orbital migration. Some calculations show that interaction with the disk can cause rapid inward migration, which, if not stopped, results in the planet reaching the "central regions still as a sub-Jovian object." More recent calculations indicate that disk evolution during migration can mitigate this problem.
Formation of stars and protoplanetary disks
Protostars
Stars are thought to form inside giant clouds of cold molecular hydrogen—giant molecular clouds roughly 300,000 times the mass of the Sun () and 20 parsecs in diameter. Over millions of years, giant molecular clouds are prone to collapse and fragmentation. These fragments then form small, dense cores, which in turn collapse into stars. The cores range in mass from a fraction to several times that of the Sun and are called protostellar (protosolar) nebulae. They possess diameters of 0.01–0.1 pc (2,000–20,000 AU) and a particle number density of roughly 10,000 to 100,000 cm−3.
The initial collapse of a solar-mass protostellar nebula takes around 100,000 years. Every nebula begins with a certain amount of angular momentum. Gas in the central part of the nebula, with relatively low angular momentum, undergoes fast compression and forms a hot hydrostatic (not contracting) core containing a small fraction of the mass of the original nebula. This core forms the seed of what will become a star. As the collapse continues, conservation of angular momentum means that the rotation of the infalling envelope accelerates, which largely prevents the gas from directly accreting onto the central core. The gas is instead forced to spread outwards near its equatorial plane, forming a disk, which in turn accretes onto the core. The core gradually grows in mass until it becomes a young hot protostar. At this stage, the protostar and its disk are heavily obscured by the infalling envelope and are not directly observable. In fact the remaining envelope's opacity is so high that even millimeter-wave radiation has trouble escaping from inside it. Such objects are observed as very bright condensations, which emit mainly millimeter-wave and submillimeter-wave radiation. They are classified as spectral Class 0 protostars. The collapse is often accompanied by bipolar outflows—jets—that emanate along the rotational axis of the inferred disk. The jets are frequently observed in star-forming regions (see Herbig–Haro (HH) objects). The luminosity of the Class 0 protostars is high — a solar-mass protostar may radiate at up to 100 solar luminosities. The source of this energy is gravitational collapse, as their cores are not yet hot enough to begin nuclear fusion.
As the infall of its material onto the disk continues, the envelope eventually becomes thin and transparent and the young stellar object (YSO) becomes observable, initially in far-infrared light and later in the visible. Around this time the protostar begins to fuse deuterium. If the protostar is sufficiently massive (above 80 Jupiter masses ()), hydrogen fusion follows. Otherwise, if its mass is too low, the object becomes a brown dwarf. This birth of a new star occurs approximately 100,000 years after the collapse begins. Objects at this stage are known as Class I protostars, which are also called young T Tauri stars, evolved protostars, or young stellar objects. By this time the forming star has already accreted much of its mass: the total mass of the disk and remaining envelope does not exceed 10–20% of the mass of the central YSO.
At the next stage the envelope completely disappears, having been gathered up by the disk, and the protostar becomes a classical T Tauri star. This happens after about 1 million years. The mass of the disk around a classical T Tauri star is about 1–3% of the stellar mass, and it is accreted at a rate of 10−7 to per year. A pair of bipolar jets is usually present as well. The accretion explains all peculiar properties of classical T Tauri stars: strong flux in the emission lines (up to 100% of the intrinsic luminosity of the star), magnetic activity, photometric variability and jets. The emission lines actually form as the accreted gas hits the "surface" of the star, which happens around its magnetic poles. The jets are byproducts of accretion: they carry away excessive angular momentum. The classical T Tauri stage lasts about 10 million years. The disk eventually disappears due to accretion onto the central star, planet formation, ejection by jets and photoevaporation by UV-radiation from the central star and nearby stars. As a result, the young star becomes a weakly lined T Tauri star, which slowly, over hundreds of millions of years, evolves into an ordinary Sun-like star.
Protoplanetary disks
Under certain circumstances the disk, which can now be called protoplanetary, may give birth to a planetary system. Protoplanetary disks have been observed around a very high fraction of stars in young star clusters. They exist from the beginning of a star's formation, but at the earliest stages are unobservable due to the opacity of the surrounding envelope. The disk of a Class 0 protostar is thought to be massive and hot. It is an accretion disk, which feeds the central protostar. The temperature can easily exceed 400 K inside 5 AU and 1,000 K inside 1 AU. The heating of the disk is primarily caused by the viscous dissipation of turbulence in it and by the infall of the gas from the nebula. The high temperature in the inner disk causes most of the volatile material—water, organics, and some rocks—to evaporate, leaving only the most refractory elements like iron. The ice can survive only in the outer part of the disk.
The main problem in the physics of accretion disks is the generation of turbulence and the mechanism responsible for the high effective viscosity. The turbulent viscosity is thought to be responsible for the transport of the mass to the central protostar and momentum to the periphery of the disk. This is vital for accretion, because the gas can be accreted by the central protostar only if it loses most of its angular momentum, which must be carried away by the small part of the gas drifting outwards. The result of this process is the growth of both the protostar and of the disk radius, which can reach 1,000 AU if the initial angular momentum of the nebula is large enough. Large disks are routinely observed in many star-forming regions such as the Orion nebula.
The lifespan of the accretion disks is about 10 million years. By the time the star reaches the classical T-Tauri stage, the disk becomes thinner and cools. Less volatile materials start to condense close to its center, forming 0.1–1 μm dust grains that contain crystalline silicates. The transport of the material from the outer disk can mix these newly formed dust grains with primordial ones, which contain organic matter and other volatiles. This mixing can explain some peculiarities in the composition of Solar System bodies such as the presence of interstellar grains in primitive meteorites and refractory inclusions in comets.
Dust particles tend to stick to each other in the dense disk environment, leading to the formation of larger particles up to several centimeters in size. The signatures of the dust processing and coagulation are observed in the infrared spectra of the young disks. Further aggregation can lead to the formation of planetesimals measuring 1 km across or larger, which are the building blocks of planets. Planetesimal formation is another unsolved problem of disk physics, as simple sticking becomes ineffective as dust particles grow larger.
One hypothesis is formation by gravitational instability. Particles several centimeters in size or larger slowly settle near the middle plane of the disk, forming a very thin—less than 100 km—and dense layer. This layer is gravitationally unstable and may fragment into numerous clumps, which in turn collapse into planetesimals. However, the differing velocities of the gas disk and the solids near the mid-plane can generate turbulence which prevents the layer from becoming thin enough to fragment due to gravitational instability. This may limit the formation of planetesimals via gravitational instabilities to specific locations in the disk where the concentration of solids is enhanced.
Another possible mechanism for the formation of planetesimals is the streaming instability in which the drag felt by particles orbiting through gas creates a feedback effect causing the growth of local concentrations. These local concentrations push back on the gas creating a region where the headwind felt by the particles is smaller. The concentration is thus able to orbit faster and undergoes less radial drift. Isolated particles join these concentrations as they are overtaken or as they drift inward causing it to grow in mass. Eventually these concentrations form massive filaments which fragment and undergo gravitational collapse forming planetesimals the size of the larger asteroids.
Planetary formation can also be triggered by gravitational instability within the disk itself, which leads to its fragmentation into clumps. Some of them, if they are dense enough, will collapse, which can lead to rapid formation of gas giant planets and even brown dwarfs on the timescale of 1,000 years. If these clumps migrate inward as the collapse proceeds tidal forces from the star can result in a significant mass loss leaving behind a smaller body. However it is only possible in massive disks—more massive than . In comparison, typical disk masses are . Because the massive disks are rare, this mechanism of planet formation is thought to be infrequent. On the other hand, it may play a major role in the formation of brown dwarfs.
The ultimate dissipation of protoplanetary disks is triggered by a number of different mechanisms. The inner part of the disk is either accreted by the star or ejected by the bipolar jets, whereas the outer part can evaporate under the star's powerful UV radiation during the T Tauri stage or by nearby stars. The gas in the central part can either be accreted or ejected by the growing planets, while the small dust particles are ejected by the radiation pressure of the central star. What is finally left is either a planetary system, a remnant disk of dust without planets, or nothing, if planetesimals failed to form.
Because planetesimals are so numerous, and spread throughout the protoplanetary disk, some survive the formation of a planetary system. Asteroids are understood to be left-over planetesimals, gradually grinding each other down into smaller and smaller bits, while comets are typically planetesimals from the farther reaches of a planetary system. Meteorites are samples of planetesimals that reach a planetary surface, and provide a great deal of information about the formation of the Solar System. Primitive-type meteorites are chunks of shattered low-mass planetesimals, where no thermal differentiation took place, while processed-type meteorites are chunks from shattered massive planetesimals. Interstellar objects could have been captured, and become part of the young Solar system.
Formation of planets
Rocky planets
According to the solar nebular disk model, rocky planets form in the inner part of the protoplanetary disk, within the frost line, where the temperature is high enough to prevent condensation of water ice and other substances into grains. This results in coagulation of purely rocky grains and later in the formation of rocky planetesimals. Such conditions are thought to exist in the inner 3–4 AU part of the disk of a Sun-like star.
After small planetesimals—about 1 km in diameter—have formed by one way or another, runaway accretion begins. It is called runaway because the mass growth rate is proportional to , where R and M are the radius and mass of the growing body, respectively. The specific (divided by mass) growth accelerates as the mass increases. This leads to the preferential growth of larger bodies at the expense of smaller ones. The runaway accretion lasts between 10,000 and 100,000 years and ends when the largest bodies exceed approximately 1,000 km in diameter. Slowing of the accretion is caused by gravitational perturbations by large bodies on the remaining planetesimals. In addition, the influence of larger bodies stops further growth of smaller bodies.
The next stage is called oligarchic accretion. It is characterized by the dominance of several hundred of the largest bodies—oligarchs, which continue to slowly accrete planetesimals. No body other than the oligarchs can grow. At this stage the rate of accretion is proportional to R2, which is derived from the geometrical cross-section of an oligarch. The specific accretion rate is proportional to ; and it declines with the mass of the body. This allows smaller oligarchs to catch up to larger ones. The oligarchs are kept at the distance of about (= is the Hill radius, where a is the semimajor axis, e is the orbital eccentricity, and Ms is the mass of the central star) from each other by the influence of the remaining planetesimals. Their orbital eccentricities and inclinations remain small. The oligarchs continue to accrete until planetesimals are exhausted in the disk around them. Sometimes nearby oligarchs merge. The final mass of an oligarch depends on the distance from the star and surface density of planetesimals and is called the isolation mass. For the rocky planets it is up to , or one Mars mass. The final result of the oligarchic stage is the formation of about 100 Moon- to Mars-sized planetary embryos uniformly spaced at about . They are thought to reside inside gaps in the disk and to be separated by rings of remaining planetesimals. This stage is thought to last a few hundred thousand years.
The last stage of rocky planet formation is the merger stage. It begins when only a small number of planetesimals remains and embryos become massive enough to perturb each other, which causes their orbits to become chaotic. During this stage embryos expel remaining planetesimals, and collide with each other. The result of this process, which lasts for 10 to 100 million years, is the formation of a limited number of Earth-sized bodies. Simulations show that the number of surviving planets is on average from 2 to 5. In the Solar System they may be represented by Earth and Venus. Formation of both planets required merging of approximately 10–20 embryos, while an equal number of them were thrown out of the Solar System. Some of the embryos, which originated in the asteroid belt, are thought to have brought water to Earth. Mars and Mercury may be regarded as remaining embryos that survived that rivalry. Rocky planets which have managed to coalesce settle eventually into more or less stable orbits, explaining why planetary systems are generally packed to the limit; or, in other words, why they always appear to be at the brink of instability.
Giant planets
The formation of giant planets is an outstanding problem in the planetary sciences. In the framework of the solar nebular model two theories for their formation exist. The first one is the disk instability model, where giant planets form in the massive protoplanetary disks as a result of its gravitational fragmentation (see above). The second possibility is the core accretion model, which is also known as the nucleated instability model. The latter scenario is thought to be the most promising one, because it can explain the formation of the giant planets in relatively low-mass disks (less than ). In this model giant planet formation is divided into two stages: a) accretion of a core of approximately and b) accretion of gas from the protoplanetary disk. Either method may also lead to the creation of brown dwarfs. Searches as of 2011 have found that core accretion is likely the dominant formation mechanism.
Giant planet core formation is thought to proceed roughly along the lines of the terrestrial planet formation. It starts with planetesimals that undergo runaway growth, followed by the slower oligarchic stage. Hypotheses do not predict a merger stage, due to the low probability of collisions between planetary embryos in the outer part of planetary systems. An additional difference is the composition of the planetesimals, which in the case of giant planets form beyond the so-called frost line and consist mainly of ice—the ice to rock ratio is about 4 to 1. This enhances the mass of planetesimals fourfold. However, the minimum mass nebula capable of terrestrial planet formation can only form cores at the distance of Jupiter (5 AU) within 10 million years. The latter number represents the average lifetime of gaseous disks around Sun-like stars. The proposed solutions include enhanced mass of the disk—a tenfold increase would suffice; protoplanet migration, which allows the embryo to accrete more planetesimals; and finally accretion enhancement due to gas drag in the gaseous envelopes of the embryos. Some combination of the above-mentioned ideas may explain the formation of the cores of gas giant planets such as Jupiter and perhaps even Saturn. The formation of planets like Uranus and Neptune is more problematic, since no theory has been capable of providing for the in situ formation of their cores at the distance of 20–30 AU from the central star. One hypothesis is that they initially accreted in the Jupiter-Saturn region, then were scattered and migrated to their present location. Another possible solution is the growth of the cores of the giant planets via pebble accretion. In pebble accretion objects between a cm and a meter in diameter falling toward a massive body are slowed enough by gas drag for them to spiral toward it and be accreted. Growth via pebble accretion may be as much as 1000 times faster than by the accretion of planetesimals.
Once the cores are of sufficient mass (), they begin to gather gas from the surrounding disk. Initially it is a slow process, increasing the core masses up to in a few million years. After that, the accretion rates increase dramatically and the remaining 90% of the mass is accumulated in approximately 10,000 years. The accretion of gas stops when the supply from the disk is exhausted. This happens gradually, due to the formation of a density gap in the protoplanetary disk and to disk dispersal. In this model ice giants—Uranus and Neptune—are failed cores that began gas accretion too late, when almost all gas had already disappeared. The post-runaway-gas-accretion stage is characterized by migration of the newly formed giant planets and continued slow gas accretion. Migration is caused by the interaction of the planet sitting in the gap with the remaining disk. It stops when the protoplanetary disk disappears or when the end of the disk is attained. The latter case corresponds to the so-called hot Jupiters, which are likely to have stopped their migration when they reached the inner hole in the protoplanetary disk.
During the accretion of gas via streams, a giant planet can be surrounded by a circumplanetary disk. This circumplanetary disk also carries solids and can form satellites. The Galilean moons are thought to have formed in such a circumplanetary disk.
Giant planets can significantly influence terrestrial planet formation. The presence of giants tends to increase eccentricities and inclinations (see Kozai mechanism) of planetesimals and embryos in the terrestrial planet region (inside 4 AU in the Solar System). If giant planets form too early, they can slow or prevent inner planet accretion. If they form near the end of the oligarchic stage, as is thought to have happened in the Solar System, they will influence the merges of planetary embryos, making them more violent. As a result, the number of terrestrial planets will decrease and they will be more massive. In addition, the size of the system will shrink, because terrestrial planets will form closer to the central star. The influence of giant planets in the Solar System, particularly that of Jupiter, is thought to have been limited because they are relatively remote from the terrestrial planets.
The region of a planetary system adjacent to the giant planets will be influenced in a different way. In such a region, eccentricities of embryos may become so large that the embryos pass close to a giant planet, which may cause them to be ejected from the system. If all embryos are removed, then no planets will form in this region. An additional consequence is that a huge number of small planetesimals will remain, because giant planets are incapable of clearing them all out without the help of embryos. The total mass of remaining planetesimals will be small, because cumulative action of the embryos before their ejection and giant planets is still strong enough to remove 99% of the small bodies. Such a region will eventually evolve into an asteroid belt, which is a full analog of the asteroid belt in the Solar System, located from 2 to 4 AU from the Sun.
Exoplanets
Thousands of exoplanets have been identified in the last twenty years, with, at the very least, billions more, within our observable universe, yet to be discovered. The orbits of many of these planets and systems of planets differ significantly from the planets in the Solar System. The exoplanets discovered include hot-Jupiters, warm-Jupiters, super-Earths, and systems of tightly packed inner planets.
The hot-Jupiters and warm-Jupiters are thought to have migrated to their current orbits during or following their formation. A number of possible mechanisms for this migration have been proposed. Type I or Type II migration could smoothly decrease the semimajor axis of the planet's orbit resulting in a warm- or hot-Jupiter. Gravitational scattering by other planets onto eccentric orbits with a perihelion near the star followed by the circularization of its orbit due to tidal interactions with the star can leave a planet on a close orbit. If a massive companion planet or star on an inclined orbit was present an exchange of inclination for eccentricity via the Kozai mechanism raising eccentricities and lowering perihelion followed by circularization can also result in a close orbit. Many of the Jupiter-sized planets have eccentric orbits which may indicate that gravitational encounters occurred between the planets, although migration while in resonance can also excite eccentricities. The in situ growth of hot Jupiters from closely orbiting super Earths has also been proposed. The cores in this hypothesis could have formed locally or at a greater distance and migrated close to the star.
Super-Earths and other closely orbiting planets are thought to have either formed in situ or ex situ, that is, to have migrated inward from their initial locations. The in situ formation of closely orbiting super-Earths would require a massive disk, the migration of planetary embryos followed by collisions and mergers, or the radial drift of small solids from farther out in the disk. The migration of the super-Earths, or the embryos that collided to form them, is likely to have been Type I due to their smaller mass. The resonant orbits of some of the exoplanet systems indicates that some migration occurred in these systems, while the spacing of the orbits in many of the other systems not in resonance indicates that an instability likely occurred in those systems after the dissipation of the gas disk. The absence of Super-Earths and closely orbiting planets in the Solar System may be due to the previous formation of Jupiter blocking their inward migration.
The amount of gas a super-Earth that formed in situ acquires may depend on when the planetary embryos merged due to giant impacts relative to the dissipation of the gas disk. If the mergers happen after the gas disk dissipates terrestrial planets can form, if in a transition disk a super-Earth with a gas envelope containing a few percent of its mass may form. If the mergers happen too early runaway gas accretion may occur leading to the formation of a gas giant. The mergers begin when the dynamical friction due to the gas disk becomes insufficient to prevent collisions, a process that will begin earlier in a higher metallicity disk. Alternatively gas accretion may be limited due to the envelopes not being in hydrostatic equilibrium, instead gas may flow through the envelope slowing its growth and delaying the onset of runaway gas accretion until the mass of the core reaches 15 Earth masses.
Meaning of accretion
Use of the term "accretion disk" for the protoplanetary disk leads to confusion over the planetary accretion process.
The protoplanetary disk is sometimes referred to as an accretion disk, because while the young T Tauri-like protostar is still contracting, gaseous material may still be falling onto it, accreting on its surface from the disk's inner edge. In an accretion disk, there is a net flux of mass from larger radii toward smaller radii.
However, that meaning should not be confused with the process of accretion forming the planets. In this context, accretion refers to the process of cooled, solidified grains of dust and ice orbiting the protostar in the protoplanetary disk, colliding and sticking together and gradually growing, up to and including the high-energy collisions between sizable planetesimals.
In addition, the giant planets probably had accretion disks of their own, in the first meaning of the word. The clouds of captured hydrogen and helium gas contracted, spun up, flattened, and deposited gas onto the surface of each giant protoplanet, while solid bodies within that disk accreted into the giant planet's regular moons.
| Physical sciences | Physical cosmology | null |
212490 | https://en.wikipedia.org/wiki/Subatomic%20particle | Subatomic particle | In physics, a subatomic particle is a particle smaller than an atom. According to the Standard Model of particle physics, a subatomic particle can be either a composite particle, which is composed of other particles (for example, a baryon, like a proton or a neutron, composed of three quarks; or a meson, composed of two quarks), or an elementary particle, which is not composed of other particles (for example, quarks; or electrons, muons, and tau particles, which are called leptons). Particle physics and nuclear physics study these particles and how they interact. Most force-carrying particles like photons or gluons are called bosons and, although they have quanta of energy, do not have rest mass or discrete diameters (other than pure energy wavelength) and are unlike the former particles that have rest mass and cannot overlap or combine which are called fermions. The W and Z bosons, however, are an exception to this rule and have relatively large rest masses at approximately 80 GeV and 90 GeV respectively.
Experiments show that light could behave like a stream of particles (called photons) as well as exhibiting wave-like properties. This led to the concept of wave–particle duality to reflect that quantum-scale behave both like particles and like waves; they are sometimes called wavicles to reflect this.
Another concept, the uncertainty principle, states that some of their properties taken together, such as their simultaneous position and momentum, cannot be measured exactly. The wave–particle duality has been shown to apply not only to photons but to more massive particles as well.
Interactions of particles in the framework of quantum field theory are understood as creation and annihilation of quanta of corresponding fundamental interactions. This blends particle physics with field theory.
Even among particle physicists, the exact definition of a particle has diverse descriptions. These professional attempts at the definition of a particle include:
A particle is a collapsed wave function
A particle is a quantum excitation of a field
A particle is an irreducible representation of the Poincaré group
A particle is an observed thing
Classification
By composition
Subatomic particles are either "elementary", i.e. not made of multiple other particles, or "composite" and made of more than one elementary particle bound together.
The elementary particles of the Standard Model are:
Six "flavors" of quarks: up, down, strange, charm, bottom, and top;
Six types of leptons: electron, electron neutrino, muon, muon neutrino, tau, tau neutrino;
Twelve gauge bosons (force carriers): the photon of electromagnetism, the three W and Z bosons of the weak force, and the eight gluons of the strong force;
The Higgs boson.
All of these have now been discovered through experiments, with the latest being the top quark (1995), tau neutrino (2000), and Higgs boson (2012).
Various extensions of the Standard Model predict the existence of an elementary graviton particle and many other elementary particles, but none have been discovered as of 2021.
Hadrons
The word hadron comes from Greek and was introduced in 1962 by Lev Okun. Nearly all composite particles contain multiple quarks (and/or antiquarks) bound together by gluons (with a few exceptions with no quarks, such as positronium and muonium). Those containing few (≤ 5) quarks (including antiquarks) are called hadrons. Due to a property known as color confinement, quarks are never found singly but always occur in hadrons containing multiple quarks. The hadrons are divided by number of quarks (including antiquarks) into the baryons containing an odd number of quarks (almost always 3), of which the proton and neutron (the two nucleons) are by far the best known; and the mesons containing an even number of quarks (almost always 2, one quark and one antiquark), of which the pions and kaons are the best known.
Except for the proton and neutron, all other hadrons are unstable and decay into other particles in microseconds or less. A proton is made of two up quarks and one down quark, while the neutron is made of two down quarks and one up quark. These commonly bind together into an atomic nucleus, e.g. a helium-4 nucleus is composed of two protons and two neutrons. Most hadrons do not live long enough to bind into nucleus-like composites; those that do (other than the proton and neutron) form exotic nuclei.
By statistics
Any subatomic particle, like any particle in the three-dimensional space that obeys the laws of quantum mechanics, can be either a boson (with integer spin) or a fermion (with odd half-integer spin).
In the Standard Model, all the elementary fermions have spin 1/2, and are divided into the quarks which carry color charge and therefore feel the strong interaction, and the leptons which do not. The elementary bosons comprise the gauge bosons (photon, W and Z, gluons) with spin 1, while the Higgs boson is the only elementary particle with spin zero.
The hypothetical graviton is required theoretically to have spin 2, but is not part of the Standard Model. Some extensions such as supersymmetry predict additional elementary particles with spin 3/2, but none have been discovered as of 2023.
Due to the laws for spin of composite particles, the baryons (3 quarks) have spin either 1/2 or 3/2 and are therefore fermions; the mesons (2 quarks) have integer spin of either 0 or 1 and are therefore bosons.
By mass
In special relativity, the energy of a particle at rest equals its mass times the speed of light squared, E = mc2. That is, mass can be expressed in terms of energy and vice versa. If a particle has a frame of reference in which it lies at rest, then it has a positive rest mass and is referred to as massive.
All composite particles are massive. Baryons (meaning "heavy") tend to have greater mass than mesons (meaning "intermediate"), which in turn tend to be heavier than leptons (meaning "lightweight"), but the heaviest lepton (the tau particle) is heavier than the two lightest flavours of baryons (nucleons). It is also certain that any particle with an electric charge is massive.
When originally defined in the 1950s, the terms baryons, mesons and leptons referred to masses; however, after the quark model became accepted in the 1970s, it was recognised that baryons are composites of three quarks, mesons are composites of one quark and one antiquark, while leptons are elementary and are defined as the elementary fermions with no color charge.
All massless particles (particles whose invariant mass is zero) are elementary. These include the photon and gluon, although the latter cannot be isolated.
By decay
Most subatomic particles are not stable. All leptons, as well as baryons decay by either the strong force or weak force (except for the proton). Protons are not known to decay, although whether they are "truly" stable is unknown, as some very important Grand Unified Theories (GUTs) actually require it. The μ and τ muons, as well as their antiparticles, decay by the weak force. Neutrinos (and antineutrinos) do not decay, but a related phenomenon of neutrino oscillations is thought to exist even in vacuums. The electron and its antiparticle, the positron, are theoretically stable due to charge conservation unless a lighter particle having magnitude of electric charge e exists (which is unlikely). Its charge is not shown yet.
Other properties
All observable subatomic particles have their electric charge an integer multiple of the elementary charge. The Standard Model's quarks have "non-integer" electric charges, namely, multiple of e, but quarks (and other combinations with non-integer electric charge) cannot be isolated due to color confinement. For baryons, mesons, and their antiparticles the constituent quarks' charges sum up to an integer multiple of e.
Through the work of Albert Einstein, Satyendra Nath Bose, Louis de Broglie, and many others, current scientific theory holds that all particles also have a wave nature. This has been verified not only for elementary particles but also for compound particles like atoms and even molecules. In fact, according to traditional formulations of non-relativistic quantum mechanics, wave–particle duality applies to all objects, even macroscopic ones; although the wave properties of macroscopic objects cannot be detected due to their small wavelengths.
Interactions between particles have been scrutinized for many centuries, and a few simple laws underpin how particles behave in collisions and interactions. The most fundamental of these are the laws of conservation of energy and conservation of momentum, which let us make calculations of particle interactions on scales of magnitude that range from stars to quarks. These are the prerequisite basics of Newtonian mechanics, a series of statements and equations in Philosophiae Naturalis Principia Mathematica, originally published in 1687.
Dividing an atom
The negatively charged electron has a mass of about of that of a hydrogen atom. The remainder of the hydrogen atom's mass comes from the positively charged proton. The atomic number of an element is the number of protons in its nucleus. Neutrons are neutral particles having a mass slightly greater than that of the proton. Different isotopes of the same element contain the same number of protons but different numbers of neutrons. The mass number of an isotope is the total number of nucleons (neutrons and protons collectively).
Chemistry concerns itself with how electron sharing binds atoms into structures such as crystals and molecules. The subatomic particles considered important in the understanding of chemistry are the electron, the proton, and the neutron. Nuclear physics deals with how protons and neutrons arrange themselves in nuclei. The study of subatomic particles, atoms and molecules, and their structure and interactions, requires quantum mechanics. Analyzing processes that change the numbers and types of particles requires quantum field theory. The study of subatomic particles per se is called particle physics. The term high-energy physics is nearly synonymous to "particle physics" since creation of particles requires high energies: it occurs only as a result of cosmic rays, or in particle accelerators. Particle phenomenology systematizes the knowledge about subatomic particles obtained from these experiments.
History
The term "subatomic particle" is largely a retronym of the 1960s, used to distinguish a large number of baryons and mesons (which comprise hadrons) from particles that are now thought to be truly elementary. Before that hadrons were usually classified as "elementary" because their composition was unknown.
A list of important discoveries follows:
| Physical sciences | Physics | null |
212523 | https://en.wikipedia.org/wiki/Pegasus%20%28constellation%29 | Pegasus (constellation) | Pegasus is a constellation in the northern sky, named after the winged horse Pegasus in Greek mythology. It was one of the 48 constellations listed by the 2nd-century astronomer Ptolemy, and is one of the 88 constellations recognised today.
With an apparent magnitude varying between 2.37 and 2.45, the brightest star in Pegasus is the orange supergiant Epsilon Pegasi, also known as Enif, which marks the horse's muzzle. Alpha (Markab), Beta (Scheat), and Gamma (Algenib), together with Alpha Andromedae (Alpheratz) form the large asterism known as the Square of Pegasus. Twelve star systems have been found to have exoplanets. 51 Pegasi was the first Sun-like star discovered to have an exoplanet companion.
Mythology
The Babylonian constellation IKU (field) had four stars of which three were later part of the Greek constellation Hippos (Pegasus). Pegasus, in Greek mythology, was a winged horse with magical powers. One myth regarding his powers says that his hooves dug out a spring, Hippocrene, which blessed those who drank its water with the ability to write poetry. Pegasus was born when Perseus cut off the head of Medusa, who was impregnated by the god Poseidon. He was born with Chrysaor from Medusa's blood. Eventually, it became the horse for Bellerophon, who was asked to kill the Chimera and succeeded with the help of Athena and Pegasus. Despite this success, after the death of his children, Bellerophon asked Pegasus to take him to Mount Olympus. Though Pegasus agreed, he plummeted back to Earth after Zeus either threw a thunderbolt at him or sent a gadfly to make Pegasus buck him off.
In ancient Persia, Pegasus was depicted by al-Sufi as a complete horse facing east, unlike most other uranographers, who had depicted Pegasus as half of a horse, rising out of the ocean. In al-Sufi's depiction, Pegasus's head is made up of the stars of Lacerta the lizard. Its right foreleg is represented by β Peg and its left foreleg is represented by η Peg, μ Peg, and λ Peg; its hind legs are marked by 9 Peg. The back is represented by π Peg and μ Cyg, and the belly is represented by ι Peg and κ Peg.
In Chinese astronomy, the modern constellation of Pegasus lies in The Black Tortoise of the north (北方玄武), where the stars were classified in several separate asterisms of stars. Epsilon and Theta Pegasi are joined with Alpha Aquarii to form Wei 危 "rooftop", with Theta forming the roof apex.
In Hindu astronomy, Pegasus is contained within the 25th nakshatra lunar mansion Purva Bhadrapada. More specifically, it represented a bedstead that was a resting place for the Moon.
For the Warrau and Arawak peoples in Guyana the stars in the Great Square, corresponding to parts of Pegasus and of Andromeda, represented a barbecue, taken up to the sky by the seven hunters of the myth of Siritjo.
Characteristics
Covering 1121 square degrees, Pegasus is the seventh-largest of the 88 constellations. Pegasus is bordered by Andromeda to the north and east, Lacerta to the north, Cygnus to the northwest, Vulpecula, Delphinus and Equuleus to the west, Aquarius to the south and Pisces to the south and east. The three-letter abbreviation for the constellation, as adopted by the IAU in 1922, is "Peg". The official constellation boundaries, as set by Belgian astronomer Eugène Delporte in 1930, are defined as a polygon of 35 segments. In the equatorial coordinate system the right ascension coordinates of these borders lie between and , while the declination coordinates are between 2.33° and 36.61°. Its position in the Northern Celestial Hemisphere means that the whole constellation is visible to observers north of 53°S.
Pegasus is dominated by a roughly square asterism, although one of the stars, Delta Pegasi or Sirrah, is now officially considered to be Alpha Andromedae, part of Andromeda, and is more usually called "Alpheratz". Traditionally, the body of the horse consists of a quadrilateral formed by the stars α Peg, β Peg, γ Peg, and α And. The front legs of the winged horse are formed by two crooked lines of stars, one leading from η Peg to κ Peg and the other from μ Peg to 1 Pegasi. Another crooked line of stars from α Peg via θ Peg to ε Peg forms the neck and head; ε is the snout.
Features
Stars
Bayer catalogued what he counted as 23 stars in the constellation, giving them the Bayer designations Alpha to Psi. He saw Pi Pegasi as one star, and was uncertain of its brightness, wavering between magnitude 4 and 5. Flamsteed labelled this star 29 Pegasi, but Bode concluded that the stars 27 and 29 Pegasi should be Pi1 and Pi2 Pegasi and that Bayer had seen them as a single star. Flamsteed added lower case letters e through to y, omitting A to D as they had been used on Bayer's chart to designate neighbouring constellations and the equator. He numbered 89 stars (now with Flamsteed designations), though 6 and 11 turned out to be stars in Aquarius. Within the constellation's borders there are 177 stars of apparent magnitude 6.5 or greater.
Epsilon Pegasi, also known as Enif, marks the horse's muzzle. The brightest star in Pegasus, is an orange supergiant of spectral type K21b that is around 12 times as massive as the Sun and is around 690 light-years distant from Earth. It is an irregular variable, its apparent magnitude varying between 2.37 and 2.45. Lying near Enif is AG Pegasi, an unusual star that brightened to magnitude 6.0 around 1885 before dimming to magnitude 9. It is composed of a red giant and white dwarf, estimated to be around 2.5 and 0.6 times the mass of the Sun respectively. With its outburst taking over 150 years, it has been described as the slowest nova ever recorded.
Three stars with Bayer designations that lie within the Great Square are variable stars. Phi and Psi Pegasi are pulsating red giants, while Tau Pegasi (the proper name is Salm), is a Delta Scuti variable—a class of short period (six hours at most) pulsating stars that have been used as standard candles and as subjects to study astroseismology. Rotating rapidly with a projected rotational velocity of 150 km s−1, Kerb is almost 30 times as luminous as the Sun and has a pulsation period of 56.5 minutes. With an outer atmosphere at an effective temperature of 7,762 K, it is a white star with a spectral type of A5IV.
Zeta, Xi, Rho and Sigma Pegasi mark the horse's neck. The brightest of these with a magnitude of 3.4 is Zeta, also traditionally known as Homam. Lying seven degrees southwest of Markab, it is a blue-white main sequence star of spectral type B8V located around 209 light-years distant. It is a slowly pulsating B star that varies slightly in luminosity with a period of 22.952 ± 0.804 hours, completing 1.04566 cycles per day. Xi lies 2 degrees northeast, and is a yellow-white main sequence star of spectral type F6V that is 86% larger and 17% more massive that the Sun, and radiate 4.5 times the solar luminosity. It has a red dwarf companion that is 192.3 au distant. If (as is likely) the smaller star is in orbit around the larger star, then it would take around 2000 years to complete a revolution. Theta Pegasi marks the horse's eye. Also known as Biham, it is a 3.43-magnitude white main sequence star of spectral type A2V, around 1.8 times as massive, 24 times as luminous, and 2.3 times as wide as the Sun.
Alpha (Markab), Beta (Scheat), and Gamma (Algenib), together with Alpha Andromedae (Alpheratz or Sirrah) form the large asterism known as the Square of Pegasus. The brightest of these, Alpheratz was also known as both Delta Pegasi and Alpha Andromedae before being placed in Andromeda in 1922 with the setting of constellation boundaries. The second brightest star is Scheat, a red giant of spectral type M2.5II-IIIe located around 196 light-years away from Earth. It has expanded until it is some 95 times as large, and has a total luminosity 1,500 times that of the Sun. Beta Pegasi is a semi-regular variable that varies from magnitude 2.31 to 2.74 over a period of 43.3 days. Markab and Algenib are blue-white stars of spectral types B9III and B2IV located 133 and 391 light-years distant respectively. Appearing to have moved off the main sequence as their core hydrogen supply is being or has been exhausted, they are enlarging and cooling to eventually become red giant stars. Markab has an apparent magnitude of 2.48, while Algenib is a Beta Cephei variable that varies between magnitudes 2.82 and 2.86 every 3 hours 38 minutes, and also exhibits some slow pulsations every 1.47 days.
Eta and Omicron Pegasi mark the left knee and Pi Pegasi the left hoof, while Iota and Kappa Pegasi mark the right knee and hoof. Also known as Matar, Eta Pegasi is the fifth-brightest star in the constellation. Shining with an apparent magnitude of 2.94, it is a multiple star system composed of a yellow giant of spectral type G2 and a yellow-white main sequence star of spectral type A5V that are 3.2 and 2.0 times as massive as the Sun. The two revolve around each other every 2.24 years. Farther afield is a binary system of two G-type main sequence stars, that would take 170,000 years to orbit the main pair if they are in fact related. Omicron Pegasi has a magnitude of 4.79. Located 300 ± 20 light-years distant from Earth, it is a white subgiant that has begun to cool, expand and brighten as it exhausts its core hydrogen fuel and moves off the main sequence. Pi1 and Pi2 Pegasi appear as an optical double to the unaided eye as they are separated by 10 arcminutes, and are not a true binary system. Located 289 ± 8 light-years distant, Pi1 is an ageing yellow giant of spectral type G6III, 1.92 times as massive and around 200 times as luminous as the Sun. Pi2 is a yellow-white subgiant that is 2.5 times as massive as the Sun and has expanded to 8 times the Sun's radius and brightened to 92 times the Sun's luminosity. It is surrounded by a circumstellar disk spinning at 145 km a second, and is 263 ± 4 light-years distant from Earth.
IK Pegasi is a close binary comprising an A-type main-sequence star and white dwarf in very close orbit; the latter a candidate for a future type Ia supernova as its main star runs out of core hydrogen fuel and expands into a giant and transfers material to the smaller star.
Twelve star systems have been found to have exoplanets. 51 Pegasi was the first Sun-like star discovered to have an exoplanet companion; 51 Pegasi b (unofficially named Bellerophon, officially named Dimidium) is a hot Jupiter close to its star, completing an orbit every four days. Spectroscopic analysis of HD 209458 b, an extrasolar planet in this constellation, has provided the first evidence of atmospheric water vapor beyond the Solar System, while extrasolar planets orbiting the star HR 8799 also in Pegasus are the first to be directly imaged. V391 Pegasi is a hot subdwarf star that has been found to have a planetary companion.
Named stars
Deep-sky objects
M15 (NGC 7078) is a globular cluster of magnitude 6.4, 34,000 light-years from Earth. It is a Shapley class IV cluster, which means that it is fairly rich and concentrated towards its center. M15 was discovered in 1746 by Jean-Dominique Maraldi. Pease 1 is a planetary nebula located within the globular cluster and was the first planetary nebula known to exist within a globular cluster. It has an apparent magnitude of 15.5.
NGC 7331 is a spiral galaxy located in Pegasus, 38 million light-years distant with a redshift of 0.0027. It was discovered by musician-astronomer William Herschel in 1784 and was later one of the first nebulous objects to be described as "spiral" by William Parsons. Another of Pegasus's galaxies is NGC 7742, a Type 2 Seyfert galaxy. Located at a distance of 77 million light-years with a redshift of 0.00555, it is an active galaxy with a supermassive black hole at its core. Its characteristic emission lines are produced by gas moving at high speeds around the central black hole.
Pegasus is also noted for its more unusual galaxies and exotic objects. Einstein's Cross is a quasar that has been lensed by a foreground galaxy. The elliptical galaxy is 400 million light-years away with a redshift of 0.0394, but the quasar is 8 billion light-years away. The lensed quasar resembles a cross because the gravitational force of the foreground galaxy on its light creates four images of the quasar. Stephan's Quintet is another unique object located in Pegasus. It is a cluster of five galaxies at a distance of 300 million light-years and a redshift of 0.0215. First discovered by Édouard Stephan, a Frenchman, in 1877, the Quintet is unique for its interacting galaxies. Two of the galaxies in the middle of the group have clearly begun to collide, sparking massive bursts of star formation and drawing off long "tails" of stars. Astronomers have predicted that all five galaxies may eventually merge into one large elliptical galaxy.
Namesakes
USS Pegasus (AK-48) and USS Pegasus (PHM-1) are United States navy ships named after the constellation "Pegasus".
The Beyblade top Storm Pegasus 105RF and its evolutions Galaxy Pegasus W105R2F and Cosmic Pegasus F:D are based on Pegasus constellation.
Pegasus Seiya, main character from the manga and anime Saint Seiya, was named after the constellation Pegasus.
| Physical sciences | Other | Astronomy |
212598 | https://en.wikipedia.org/wiki/Minor-planet%20moon | Minor-planet moon | A minor-planet moon is an astronomical object that orbits a minor planet as its natural satellite. , there are 457 minor planets known or suspected to have moons. Discoveries of minor-planet moons (and binary objects, in general) are important because the determination of their orbits provides estimates on the mass and density of the primary, allowing insights into their physical properties that are generally not otherwise accessible.
Several of the moons are quite large compared to their primaries: 90 Antiope, Mors–Somnus and Sila–Nunam (95%), Patroclus–Menoetius, Altjira and Lempo–Hiisi (90%, with Lempo–Paha at 50%). The largest known minor-planet moon in absolute size is Pluto's largest moon Charon, which itself has about half the diameter of Pluto.
There are also several known ring systems around distant objects (see: Rings of Chariklo and Chiron).
Terminology
In addition to the terms satellite and moon, the term "binary" (binary minor planet) is sometimes used for minor planets with one moon, and "triple" for minor planets with two moons. If one object is much bigger it is referred to as the primary and its companion as the secondary. The term double asteroid is sometimes used for systems in which the asteroid and its moon are roughly the same size, while binary tends to be used independently from the relative sizes of the components. When binary minor planets are similar in size, the Minor Planet Center (MPC) refers to them as "binary companions" instead of referring to the smaller body as a satellite. A good example of a true binary is the 90 Antiope system, identified in August 2000. Very small satellites are often referred to as moonlets.
Discovery milestones
Prior to the era of the Hubble Space Telescope and space probes reaching the outer Solar System, attempts to detect satellites around asteroids were limited to optical observations from Earth. For example, in 1978, stellar occultation observations were claimed as evidence of a satellite for the asteroid 532 Herculina. However, later more-detailed imaging by the Hubble Telescope did not reveal a satellite, and the current consensus is that Herculina does not have a significant satellite. There were other similar reports of asteroids having companions (usually referred to as satellites) in the following years. A letter by astronomer Thomas Hamilton in the Sky & Telescope magazine at this time pointed to apparently simultaneous impact craters on Earth (for example, the Clearwater Lakes in Quebec), suggesting that these craters were caused by pairs of gravitationally bound objects.
Also in 1978, Pluto's largest moon Charon was discovered; however, at the time Pluto was still considered to be one of the major planets.
In 1993, the first asteroid moon was confirmed when the Galileo probe discovered the small Dactyl orbiting 243 Ida in the asteroid belt. The second was discovered around 45 Eugenia in 1998. In 2001, 617 Patroclus and its same-sized companion Menoetius became the first known binary asteroids in the Jupiter trojans. The first trans-Neptunian binary after Pluto–Charon, , was optically resolved in 2002.
Multiple systems
In 2005, the asteroid 87 Sylvia was discovered to have two satellites, making it the first known triple system (also called triple minor planets). This was followed by the discovery of a second moon orbiting 45 Eugenia. Also in 2005, the dwarf planet was discovered to have two moons, making it the second trans-Neptunian object after Pluto known to have more than one moon. Additionally, 216 Kleopatra and 93 Minerva were discovered to be triple asteroids in 2008 and 2009 respectively. Since the first few triple minor planets were discovered, more continue to be discovered. , the total number of known multiple systems among minor planets is 18 (including the Pluto and Haumea systems).
The following table lists all satellites of multiple systems, starting with Pluto, which was unnumbered when its first moon was discovered in 1978. The highest known multiplicities are for Pluto (a sextuple system) and 130 Elektra (a quadruple system).
Commonality
The data about the populations of binary objects are still patchy. In addition to the inevitable observational bias (dependence on the distance from Earth, size, albedo and separation of the components) the frequency appears to be different among different categories of objects. Among asteroids, an estimated 2% would have satellites. Among trans-Neptunian objects (TNOs), an estimated 11% are thought to be binary or multiple objects, and the majority of the large TNOs have at least one satellite, including all four IAU-listed dwarf planets.
More than 50 binaries are known in each of the main groupings: near-Earth asteroids, belt asteroids, and trans-Neptunian objects, not including numerous claims based solely on light-curve variation.
Two binaries have been found so far among centaurs with semi-major axes smaller than Neptune. Both are double ring systems around 2060 Chiron and 10199 Chariklo, discovered in 1993–2011 and 2013 respectively.
Origin
The origin of minor-planet moons is not currently known with certainty, and a variety of hypotheses exist. One such model is that minor-planet moons are formed from debris knocked off the primary by an impact. Other pairings may be formed when a small object is captured by the gravity of a larger one.
Formation by collision is constrained by the angular momentum of the components, i.e. by the masses and their separation. Close binaries fit this model (e.g. Pluto–Charon). Distant binaries however, with components of comparable size, are unlikely to have followed this scenario, unless considerable mass has been lost in the event.
The distances of the components for the known binaries vary from a few hundreds of kilometres (243 Ida, 3749 Balam) to more than 3000 km (379 Huenna) for the asteroids. Among TNOs, the known separations vary from 3,000 to 50,000 km.
Populations and classes
What is "typical" for a binary system tends to depend on its location in the Solar System (presumably because of different modes of origin and lifetimes of such systems in different populations of minor planets).
Among near-Earth asteroids, satellites tend to orbit at distances of the order of 3–7 primary radii, and have diameters two to several times smaller than the primary. Since these binaries are all inner-planet crossers, it is thought that tidal stresses that occurred when the parent object passed close to a planet may be responsible for the formation of many of them, although collisions are thought to also be a factor in the creation of these satellites.
Among main-belt asteroids, the satellites are usually much smaller than the primary (a notable exception being 90 Antiope), and orbit around 10 primary radii away. Many of the binary systems here are members of asteroid families, and a good proportion of satellites are expected to be fragments of a parent body whose disruption after an asteroid collision produced both the primary and satellite.
Among trans-Neptunian objects, it is common for the two orbiting components to be of comparable size, and for the semi-major axis of their orbits to be much larger − about 100 to 1000 primary radii. A significant proportion of these binaries are expected to be primordial.
Pluto has five known moons. Its largest moon Charon has a radius of more than half that of Pluto, and is large enough to orbit a point outside Pluto's surface. In fact, each orbits the common barycenter between them, with Pluto's orbit entirely enclosed by Charon's; thus they form a binary system informally referred to as a double dwarf planet. Pluto's four other moons, Nix, Hydra, Kerberos, and Styx, are far smaller and orbit the Pluto–Charon system.
Haumea has two moons with radii estimated around 155 km (Hiʻiaka) and 85 km (Namaka).
has one known moon, S/2015 (136472) 1, estimated to be some in diameter.
47171 Lempo is a unique trans-Neptunian triple system: Lempo and its moon of roughly equal mass, Hiisi, form a close-proximity binary, separated by roughly 867 km. A second moon, Paha, orbits the Lempo–Hiisi binary at about 7411 km.
has one known moon, Dysnomia. Its radius, based on its brightness, is estimated to be roughly between 150 and 350 km.
List
, there are 457 minor planets (systems) with 477 known companions. The following table is a listing of the total number of these systems by orbital class:
Near-Earth objects
This is a list of near-Earth asteroids with companions. Candidate binaries with an unconfirmed status are displayed on a dark background. For an overview, see summary and introduction.
Mars crossers
This is a list of Mars-crossing asteroids with companions. Candidate binaries with an unconfirmed status are displayed on a dark background. For an overview, see summary and introduction.
Main-belt asteroids
This is a list of main-belt asteroids with companions. Candidate binaries with an unconfirmed status are displayed on a dark background. For an overview, see summary and introduction.
The following binaries are double asteroids, with similarly sized components, and a barycenter outside of the larger objec
90 Antiope – S/2000 (90) 1
854 Frostia – S/2004 (854) 1
1313 Berna – S/2004 (1313) 1
2478 Tokai – S/2007 (2478) 1
3169 Ostro – S/2005 (3169) 1
3749 Balam – S/2002 (3749) 1
3905 Doppler – S/2013 (3905) 1
4674 Pauling – S/2004 (4674) 1
4951 Iwamoto – S/2007 (4951) 1
5674 Wolff – S/2015 (5674) 1
8474 Rettig – S/2015 (8474) 1
17246 Christophedumas – S/2004 (17246) 1
– S/2011 (300163) 1
In addition, these bodies might be double asteroids, but due to errors in their size and orbit, it is uncertain.
809 Lundia – S/2005 (809) 1
1089 Tama – S/2003 (1089) 1
1509 Esclangona – S/2003 (1509) 1
4492 Debussy – S/2004 (4492) 1
11264 Claudiomaccone – S/2003 (11264) 1
22899 Alconrad – S/2003 (22899) 1
Jupiter trojans
This is a list of Jupiter trojans with companions. Candidate binaries with an unconfirmed status are displayed on a dark background. For an overview, see summary and introduction.
Trans-Neptunian objects
This is a list of trans-Neptunian objects with companions. Candidate binaries with an unconfirmed status are displayed on a dark background. This list gives the companion's orbital period (Ps) in days rather than hours. For an overview, see summary and introduction.
| Physical sciences | Planetary science | Astronomy |
212698 | https://en.wikipedia.org/wiki/Quackery | Quackery | Quackery, often synonymous with health fraud, is the promotion of fraudulent or ignorant medical practices. A quack is a "fraudulent or ignorant pretender to medical skill" or "a person who pretends, professionally or publicly, to have skill, knowledge, qualification or credentials they do not possess; a charlatan or snake oil salesman". The term quack is a clipped form of the archaic term , derived from a "hawker of salve" or rather somebody who boasted about their salves, more commonly known as ointments. In the Middle Ages the term quack meant "shouting". The quacksalvers sold their wares at markets by shouting to gain attention.
Common elements of general quackery include questionable diagnoses using questionable diagnostic tests, as well as untested or refuted treatments, especially for serious diseases such as cancer. Quackery is often described as "health fraud" with the salient characteristic of aggressive promotion.
Definition
Psychiatrist and author Stephen Barrett of Quackwatch defines quackery as "the promotion of unsubstantiated methods that lack a scientifically plausible rationale" and more broadly as:
In addition to the ethical problems of promising benefits that are not likely to occur, quackery might cause people to forego treatments that are more likely to help them, in favor of ineffective treatments given by the "quack".
American pediatrician Paul Offit has proposed four ways in which alternative medicine "becomes quackery":
"by recommending against conventional therapies that are helpful."
"by promoting potentially harmful therapies without adequate warning."
"by draining patients' bank accounts ..."
"by promoting magical thinking ..."
Since it is difficult to distinguish between those who knowingly promote unproven medical therapies and those who are mistaken as to their effectiveness, United States courts have ruled in defamation cases that accusing someone of quackery or calling a practitioner a quack is not equivalent to accusing that person of committing medical fraud. However, the FDA makes little distinction between the two. To be considered a fraud, it is not strictly necessary for one to know they are misrepresenting the benefits or risks of the services offered.
Quacksalver
Unproven, usually ineffective, and sometimes dangerous medicines and treatments have been peddled throughout human history. Theatrical performances were sometimes given to enhance the credibility of purported medicines. Grandiose claims were made for what could be humble materials indeed: for example, in the mid-19th century revalenta arabica was advertised as having extraordinary restorative virtues as an empirical diet for invalids; despite its impressive name and many glowing testimonials it was in truth only ordinary lentil flour, sold to the gullible at many times the true cost.
Even where no fraud was intended, quack remedies often contained no effective ingredients whatsoever. Some remedies contained substances such as opium, alcohol and honey, which would have given symptomatic relief but had no curative properties. Some would have addictive qualities to entice the buyer to return. The few effective remedies sold by quacks included emetics, laxatives and diuretics. Some ingredients did have medicinal effects: mercury, silver and arsenic compounds may have helped some infections and infestations; willow bark contains salicylic acid, chemically closely related to aspirin; and the quinine contained in Jesuit's bark was an effective treatment for malaria and other fevers. However, knowledge of appropriate uses and dosages was limited.
Criticism of quackery in academia
The evidence-based medicine community has criticized the infiltration of alternative medicine into mainstream academic medicine, education, and publications, accusing institutions of "diverting research time, money, and other resources from more fruitful lines of investigation in order to pursue a theory that has no basis in biology."
For example, David Gorski criticized Brian M. Berman, founder of the University of Maryland Center for Integrative Medicine, for writing that "There [is] evidence that both real acupuncture and sham acupuncture [are] more effective than no treatment and that acupuncture can be a useful supplement to other forms of conventional therapy for low back pain." He also castigated editors and peer reviewers at the New England Journal of Medicine for allowing it to be published, since it effectively recommended deliberately misleading patients in order to achieve a known placebo effect.
History in Europe and the United States
With little understanding of the causes and mechanisms of illnesses, widely marketed "cures" (as opposed to locally produced and locally used remedies), often referred to as patent medicines, first came to prominence during the 17th and 18th centuries in Britain and the British colonies, including those in North America. Daffy's Elixir and Turlington's Balsam were among the first products that used branding (e.g. using highly distinctive containers) and mass marketing to create and maintain markets. A similar process occurred in other countries of Europe around the same time, for example with the marketing of Eau de Cologne as a cure-all medicine by Johann Maria Farina and his imitators. Patent medicines often contained alcohol or opium, which, while presumably not curing the diseases for which they were sold as a remedy, did make the imbibers feel better and confusedly appreciative of the product.
The number of internationally marketed quack medicines increased in the later 18th century; the majority of them originated in Britain and were exported throughout the British Empire. By 1830, British parliamentary records list over 1,300 different "proprietary medicines", the majority of which were "quack" cures by modern standards.
A Dutch organisation that opposes quackery, (VtdK), was founded in 1881, making it the oldest organisation of this kind in the world. It has published its magazine (Dutch Magazine against Quackery) ever since. In these early years the played a part in the professionalisation of medicine. Its efforts in the public debate helped to make the Netherlands one of the first countries with governmental drug regulation.
In 1909, in an attempt to stop the sale of quack medicines, the British Medical Association published Secret Remedies, What They Cost And What They Contain. This publication was originally a series of articles published in the British Medical Journal between 1904 and 1909. The publication was composed of 20 chapters, organising the work by sections according to the ailments the medicines claimed to treat. Each remedy was tested thoroughly, the preface stated: "Of the accuracy of the analytical data there can be no question; the investigation has been carried out with great care by a skilled analytical chemist." The book did lead to the end of some of the quack cures, but some survived the book by several decades. For example, Beecham's Pills, which according to the British Medical Association contained in 1909 only aloes, ginger and soap, but claimed to cure 31 medical conditions, were sold until 1998. The failure of the medical establishment to stop quackery was rooted in the difficulty of defining what precisely distinguished real medicine, and in the appeals that quackery held out to consumers.
British patent medicines lost their dominance in the United States when they were denied access to the Thirteen Colonies markets during the American Revolution, and lost further ground for the same reason during the War of 1812. From the early 19th century "home-grown" American brands started to fill the gap, reaching their peak in the years after the American Civil War. British medicines never regained their previous dominance in North America, and the subsequent era of mass marketing of American patent medicines is usually considered to have been a "golden age" of quackery in the United States. This was mirrored by similar growth in marketing of quack medicines elsewhere in the world.
In the United States, false medicines in this era were often denoted by the slang term snake oil, a reference to sales pitches for the false medicines that claimed exotic ingredients provided the supposed benefits. Those who sold them were called "snake oil salesmen", and usually sold their medicines with a fervent pitch similar to a fire and brimstone religious sermon. They often accompanied other theatrical and entertainment productions that traveled as a road show from town to town, leaving quickly before the falseness of their medicine was discovered. Not all quacks were restricted to such small-time businesses however, and a number, especially in the United States, became enormously wealthy through national and international sales of their products.
In 1875, the Pacific Medical and Surgical Journal complained:
One among many examples is William Radam, a German immigrant to the US, who, in the 1880s, started to sell his "Microbe Killer" throughout the United States and, soon afterwards, in Britain and throughout the British colonies. His concoction was widely advertised as being able to "cure all diseases", and this phrase was even embossed on the glass bottles the medicine was sold in. In fact, Radam's medicine was a therapeutically useless (and in large quantities actively poisonous) dilute solution of sulfuric acid, coloured with a little red wine. Radam's publicity material, particularly his books, provide an insight into the role that pseudoscience played in the development and marketing of "quack" medicines towards the end of the 19th century.
Advertising claims similar to those of Radam can be found throughout the 18th, 19th, 20th and 21st centuries. "Dr." Sibley, an English patent medicine seller of the late 18th and early 19th centuries, even went so far as to claim that his Reanimating Solar Tincture would, as the name implies, "restore life in the event of sudden death". Another English quack, "Dr. Solomon" claimed that his Cordial Balm of Gilead cured almost anything, but was particularly effective against all venereal complaints, from gonorrhea to onanism. Although it was basically just brandy flavoured with herbs, the price of a bottle was a half guinea (£sd system) in 1800, equivalent to over £ ($) in 2014.
Not all patent medicines were without merit. Turlingtons Balsam of Life, first marketed in the mid-18th century, did have genuinely beneficial properties. This medicine continued to be sold under the original name into the early 20th century, and can still be found in the British and American pharmacopoeias as "Compound tincture of benzoin". In these cases, the treatments likely lacked empirical support when they were introduced to the market, and their benefits were simply a convenient coincidence discovered after the fact.
The end of the road for the quack medicines now considered grossly fraudulent in the nations of North America and Europe came in the early 20th century. 21 February 1906 saw the passage into law of the Pure Food and Drug Act in the United States. This was the result of decades of campaigning by both government departments and the medical establishment, supported by a number of publishers and journalists (one of the most effective was Samuel Hopkins Adams, who wrote "The Great American Fraud" series in Collier's in 1905). This American Act was followed three years later by similar legislation in Britain and in other European nations. Between them, these laws began to remove the more outrageously dangerous contents from patent and proprietary medicines, and to force quack medicine proprietors to stop making some of their more blatantly dishonest claims. The Act, however, left advertising and claims of effectiveness unregulated as the Supreme Court interpreted it to mean only that ingredients on labels had to be accurate. Language in the 1912 Sherley Amendment, meant to close this loophole, was limited to regulating claims that were false and fraudulent, creating the need to show intent. Throughout the early 20th century, the American Medical Association collected material on medical quackery, and one of their members and medical editors in particular, Arthur J. Cramp, devoted his career to criticizing such products. The AMA's Department of Investigation closed in 1975, but their only archive open to non-members remains, the American Medical Association Health Fraud and Alternative Medicine Collection.
"Medical quackery and promotion of nostrums and worthless drugs were among the most prominent abuses that led to formal self-regulation in business and, in turn, to the creation of the Better Business Bureau."
Contemporary culture
"Quackery is the promotion of false and unproven health schemes for a profit. It is rooted in the traditions of the marketplace", with "commercialism overwhelming professionalism in the marketing of alternative medicine". Quackery is most often used to denote the peddling of the "cure-alls" described above. Quackery is an ongoing problem that can be found in any culture and in every medical tradition. Unlike other advertising mediums, rapid advancements in communication through the Internet have opened doors for an unregulated market of quack cures and marketing campaigns rivaling the early 20th century. Most people with an e-mail account have experienced the marketing tactics of spamming – in which modern forms of quackery are touted as miraculous remedies for "weight loss" and "sexual enhancement", as well as outlets for medicines of unknown quality.
India
In 2008, the Hindustan Times reported that some officials and doctors estimated that there were more than 40,000 quacks practicing in Delhi, following outrage over a "multi-state racket where unqualified doctors conducted hundreds of illegal kidney transplants for huge profits." The president of the Indian Medical Association (IMA) in 2008 criticized the central government for failing to address the problem of quackery and for not framing any laws against it.
In 2017, IMA again asked for an antiquackery law with stringent action against those practicing without a license. As of 2024, the government of India is yet to pass an anti-quackery law.
Ministry of Ayush
In 2014, the Government of India formed a Ministry of AYUSH that includes the seven traditional systems of healthcare.
The Ministry of Ayush (expanded from Ayurveda, Yoga, Naturopathy, Unani, Siddha, Sowa-Rigpa and Homoeopathy), is purposed with developing education, research and propagation of indigenous alternative medicine systems in India. The ministry has faced significant criticism for funding systems that lack biological plausibility and are either untested or conclusively proven as ineffective. Quality of research has been poor, and drugs have been launched without any rigorous pharmacological studies and meaningful clinical trials on Ayurveda or other alternative healthcare systems.
There is no credible efficacy or scientific basis of any of these forms of treatment. A strong consensus prevails among the scientific community that homeopathy is a pseudo-scientific, unethical and implausible line of treatment. Ayurveda is deemed to be pseudoscientific. Much of the research on postural yoga has taken the form of preliminary studies or clinical trials of low methodological quality; there is no conclusive therapeutic effect except in back pain. Naturopathy is considered to be a form of pseudo-scientific quackery, ineffective and possibly harmful, with a plethora of ethical concerns about the very practice.
Unani lacks biological plausibility and is considered to be pseudo-scientific quackery, as well.
United States
While quackery is often aimed at the aged or chronically ill, it can be aimed at all age groups, including teens, and the FDA has mentioned some areas where potential quackery may be a problem: breast developers, weight loss, steroids and growth hormones, tanning and tanning pills, hair removal and growth, and look-alike drugs.
In 1992, the president of The National Council Against Health Fraud, William T. Jarvis, wrote in Clinical Chemistry that:
For those in the practice of any medicine, to allege quackery is to level a serious objection to a particular form of practice. Most developed countries have a governmental agency, such as the Food and Drug Administration (FDA) in the US, whose purpose is to monitor and regulate the safety of medications as well as the claims made by the manufacturers of new and existing products, including drugs and nutritional supplements or vitamins. The Federal Trade Commission (FTC) participates in some of these efforts. To better address less regulated products, in 2000, US President Clinton signed Executive Order 13147 that created the White House Commission on Complementary and Alternative Medicine. In 2002, the commission's final report made several suggestions regarding education, research, implementation, and reimbursement as ways to evaluate the risks and benefits of each. As a direct result, more public dollars have been allocated for research into some of these methods.
Individuals and non-governmental agencies are active in attempts to expose quackery. According to John C. Norcross et al. less is consensus about ineffective "compared to effective procedures" but identifying both "pseudoscientific, unvalidated, or 'quack' psychotherapies" and "assessment measures of questionable validity on psycho-metric grounds" was pursued by various authors. The evidence-based practice (EBP) movement in mental health emphasizes the consensus in psychology that psychological practice should rely on empirical research. There are also "anti-quackery" websites, such as Quackwatch, that help consumers evaluate claims. Quackwatch's information is relevant to both consumers and medical professionals.
Presence and acceptance
There have been several suggested reasons why quackery is accepted by patients in spite of its lack of effectiveness:
Ignorance Those who perpetuate quackery may do so to take advantage of ignorance about conventional medical treatments versus alternative treatments, or may themselves be ignorant regarding their own claims. Mainstream medicine has produced many remarkable advances, so people may tend to also believe groundless claims.
Placebo effect Medicines or treatments known to have no pharmacological effect on a disease can still affect a person's perception of their illness, and this belief in its turn does indeed sometimes have a therapeutic effect, causing the patient's condition to improve. This is not to say that no real cure of biological illness is effected – "though we might describe a placebo effect as being 'all in the mind', we now know that there is a genuine neurobiological basis to this phenomenon." People report reduced pain, increased well-being, improvement, or even total alleviation of symptoms. For some, the presence of a caring practitioner and the dispensation of medicine is curative in itself.
Regression fallacy Lack of understanding that health conditions change with no treatment and attributing changes in ailments to a given therapy.
Confirmation bias The tendency to search for, interpret, or prioritize information in a way that confirms one's beliefs or hypotheses. It is a type of cognitive bias and a systematic error of inductive reasoning.
Distrust of conventional medicine Many people, for various reasons, have a distrust of conventional medicine, or of the regulating organizations such as the FDA, or the major drug corporations. For example, "CAM may represent a response to disenfranchisement [discrimination] in conventional medical settings and resulting distrust".
Conspiracy theories Anti-quackery activists ("quackbusters") are often falsely accused of being part of a huge "conspiracy" to suppress "unconventional" and/or "natural" therapies, as well as those who promote them. It is alleged that this conspiracy is backed and funded by the pharmaceutical industry and the established medical care system – represented by the AMA, FDA, ADA, CDC, WHO, etc. – for the purpose of preserving their power and increasing their profits. This idea is often held by people with antiscience views.
Fear of side effects A great variety of pharmaceutical medications can have very distressing side effects, and many people fear surgery and its consequences, so they may opt to shy away from these mainstream treatments.
Cost There are some people who simply cannot afford conventional treatment, and seek out a cheaper alternative. Nonconventional practitioners can often dispense treatment at a much lower cost. This is compounded by reduced access to healthcare.
Desperation People with a serious or terminal disease, or who have been told by their practitioner that their condition is "untreatable", may react by seeking out treatment, disregarding the lack of scientific proof for its effectiveness, or even the existence of evidence that the method is ineffective or even dangerous. Despair may be exacerbated by the lack of palliative non-curative end-of-life care. Between 2012 and 2018 appeals on UK crowdfunding sites for cancer treatment with an alternative health element have raised £8 million. This is described as "a new and lucrative revenue stream for cranks, charlatans, and conmen who prey on the vulnerable."
Pride Once people have endorsed or defended a cure, or invested time and money in it, they may be reluctant or embarrassed to admit its ineffectiveness and therefore recommend a treatment that does not work. This is a manifestation of the sunk cost fallacy.
Fraud Some practitioners, fully aware of the ineffectiveness of their medicine, may intentionally produce fraudulent scientific studies, for example, thereby confusing any potential consumers as to the effectiveness of the medical treatment.
Deceased persons accused of quackery
Thomas Allinson (1858–1918), founder of naturopathy. His views often brought him into conflict with the Royal College of Physicians of Edinburgh and the General Medical Council, particularly his opposition to doctors' frequent use of drugs, his opposition to vaccination and his self-promotion in the press. His views and publication of them led to him being labeled a quack and being struck off by the General Medical Council for infamous conduct in a professional respect.
Lovisa Åhrberg (1801–1881), the first Swedish female doctor. Åhrberg was met with strong resistance from male doctors and was accused of quackery. During the formal examination she was acquitted of all charges and allowed to practice medicine in Stockholm even though it was forbidden for women in the 1820s. She later received a medal for her work.
Johanna Brandt (1876–1964), a South African naturopath who advocated the "Grape Cure" as a cure for cancer.
John R. Brinkley (1885–1942), a nonphysician and xenotransplant specialist in Kansas, US, who claimed to have discovered a method of effectively transplanting the testicles of goats into aging men. After state authorities took steps to shut down his practice, he retaliated by entering politics in 1930 and unsuccessfully running for the office of Governor of Kansas.
Hulda Regehr Clark (1928–2009), was a controversial naturopath, author, and practitioner of alternative medicine who claimed to be able to cure all diseases and advocated methods that have no scientific validity.
Max Gerson (1881–1959), was a German-born American physician who developed a dietary-based alternative cancer treatment that he claimed could cure cancer and most chronic, degenerative diseases. His treatment was called The Gerson Therapy. Most notably, Gerson Therapy was used, unsuccessfully, to treat Jessica Ainscough and Garry Winogrand. According to Quackwatch, Gerson Institute claims of cure are based not on actual documentation of survival, but on "a combination of the doctor's estimate that the departing patient has a 'reasonable chance of surviving', plus feelings that the Institute staff have about the status of people who call in". The American Cancer Society reports that "[t]here is no reliable scientific evidence that Gerson therapy is effective ..."
Samuel Hahnemann (1755–1843), founder of homeopathy. Hahnemann believed that all diseases were caused by "miasms", which he defined as irregularities in the patient's vital force. He also said that illnesses could be treated by substances that in a healthy person produced similar symptoms to the illness, in extremely low concentrations, with the therapeutic effect increasing with dilution and repeated shaking.
Lawrence B. Hamlin (in 1916), was fined under the 1906 US Pure Food and Drug Act for advertising that his Wizard Oil could kill cancer.
L. Ron Hubbard (1911–1986), was the founder of the Church of Scientology. He was an American science fiction writer, former US Navy officer, and creator of Dianetics. He has been commonly called a quack and a con man by both critics of Scientology and by many psychiatric organizations in part for his often extreme anti-psychiatric beliefs and false claims about technologies such as the E-meter.
Linda Hazzard, (1867–1938), was a self-declared doctor and fasting specialist, which she advertised as a panacea for every medical ailment. Up to 40 patients may have died of starvation in her "sanitarium" in Olalla, Washington, US. Imprisoned for one death in 1912, Hazzard was paroled in 1915 and continued to practice medicine without a license in New Zealand (1915–1920) and Washington, US (1920–1935). Died in 1938 while attempting a fasting to cure herself.
William Donald Kelley, (1925–2005), was an orthodontist and a follower of Max Gerson who developed his own alternative cancer treatment called Nonspecific Metabolic Therapy. This treatment is based on the unsubstantiated belief that "wrong foods [cause] malignancy to grow, while proper foods [allow] natural body defenses to work". It involves, specifically, treatment with pancreatic enzymes, 50 daily vitamins and minerals (including laetrile), frequent body shampoos, coffee enemas, and a specific diet. According to Quackwatch, "not only is his therapy ineffective, but people with cancer who take it die more quickly and have a worse quality of life than those having standard treatment, and can develop serious or fatal side-effects. Kelley's most famous patient was actor Steve McQueen.
John Harvey Kellogg (1852–1943), was a medical doctor in Battle Creek, Michigan, US, who ran a sanitarium using holistic methods, with a particular focus on nutrition, enemas and exercise. Kellogg was an advocate of vegetarianism and invented the corn flake breakfast cereal with his brother, Will Keith Kellogg.
John St. John Long (1798–1834) was an Irish artist who claimed to be able to cure tuberculosis by causing a sore or wound on the back of the patient, out of which the disease would exit. He was tried twice for manslaughter of his patients who died under this treatment.
Franz Anton Mesmer (1734–1815), was a German physician and astrologist, who invented what he called magnétisme animal.
Theodor Morell (1886–1948), a German physician best known as Adolf Hitler's personal doctor. Morell administered approximately 74 substances, in 28 different mixtures to Hitler, including heroin, cocaine, Doktor Koster's Antigaspills, potassium bromide, papaverine, testosterone, vitamins and animal enzymes. Despite Hitler's dependence on Morell, and his recommendations of him to other Nazi leaders, Hermann Göring, Heinrich Himmler, Albert Speer and others quietly dismissed Morell as a quack.
Daniel David Palmer (1845–1913), was a grocery store owner that claimed to have healed a janitor of deafness after adjusting the alignment of his back. He founded the field of chiropractic based on the principle that all disease and ailments could be fixed by adjusting the alignment of someone's back. His hypothesis was disregarded by medical professionals at the time and despite a considerable following has yet to be scientifically proven. Palmer established a magnetic healing facility in Davenport, Iowa, styling himself 'doctor'. Not everyone was convinced, as a local paper in 1894 wrote about him: "A crank on magnetism has a crazy notion that he can cure the sick and crippled with his magnetic hands. His victims are the weak-minded, ignorant and superstitious, those foolish people who have been sick for years and have become tired of the regular physician and want health by the short-cut method … he has certainly profited by the ignorance of his victims … His increase in business shows what can be done in Davenport, even by a quack."
Louis Pasteur (1822–1895), was a French chemist best known for his remarkable breakthroughs in microbiology. His experiments confirmed the germ theory of disease, also reducing mortality from puerperal fever (childbed), and he created the first vaccine for rabies. He is best known to the general public for showing how to stop milk and wine from going sourthis process came to be called pasteurization. His hypotheses initially met with much hostility, and he was accused of quackery on multiple occasions. However, he is now regarded as one of the three main founders of microbiology, together with Ferdinand Cohn and Robert Koch.
Linus Pauling (1901–1994), a Nobel Prize winner in chemistry, Pauling spent much of his later career arguing for the treatment of somatic and psychological diseases with orthomolecular medicine. Among his claims were that the common cold could be cured with massive doses of vitamin C. Together with Ewan Cameron he wrote the 1979 book Cancer and Vitamin C, which was again more popular with the public than the medical profession, which continued to regard claims about the effectiveness of vitamin C in treating or preventing cancer as quackery. A biographer has discussed how controversial his views on megadoses of Vitamin C have been and that he was "still being called a 'fraud' and a 'quack' by opponents of his 'orthomolecular medicine.
Doctor John Henry Pinkard (1866–1934) was a Roanoke, Virginia businessman and "Yarb Doctor" or "Herb Doctor" who concocted quack medicines that he sold and distributed in violation of the Food and Drugs Act and the earlier Pure Food and Drug Act. He was also known as a "clairvoyant, herb doctor and spiritualist." Some of Pinkard's Sanguinaria Compound, made from bloodroot or bloodwort, was seized by federal officials in 1931. "Analysis by this department of a sample of the article showed that it consisted essentially of extracts of plant drugs including sanguinaria, sugar, alcohol, and water. It was alleged in the information that the article was misbranded in that certain statements, designs, and devices regarding the therapeutic and curative effects of the article, appearing on the bottle label, falsely and fraudulently represented that it would be effective as a treatment, remedy, and cure for pneumonia, coughs, weak lungs, asthma, kidney, liver, bladder, or any stomach troubles, and effective as a great blood and nerve tonic." He pleaded guilty and was fined.
Wilhelm Reich (1897–1957), Austrian-American Psychoanalyst. Claimed that he had discovered a primordial cosmic energy called Orgone. He developed several devices, including the Cloudbuster and the Orgone Accumulator, that he believed could use orgone to manipulate the weather, battle space aliens and cure diseases, including cancer. After an investigation, the US Food and Drug Administration concluded that they were dealing with a "fraud of the first magnitude". On 10 February 1954, the US Attorney for Maine filed a complaint seeking a permanent injunction under Sections 301 and 302 of the Federal Food, Drug, and Cosmetic Act, to prevent interstate shipment of orgone accumulators and to ban some of Reich's writing promoting and advertising the devices. Reich refused to appear in court, arguing that no court was in a position to evaluate his work. Reich was arrested for contempt of court, and convicted to two years in jail, a US$10,000 fine, and his Orgone Accumulators and work on Orgone were ordered to be destroyed. On 23 August 1956, six tons of his books, journals, and papers were burned in the 25th Street public incinerator in New York. On 12 March 1957 he was sent to Danbury Federal Prison, where Richard C. Hubbard, a psychiatrist who admired Reich, examined him, recording paranoia manifested by delusions of grandiosity, persecution, and ideas of reference. Nine months later, on 18 November 1957, Reich died of a heart attack while he was in the federal penitentiary in Lewisburg, Pennsylvania.
William Herbert Sheldon (1898–1977), who created the theory of somatotypes corresponding to intelligence.
Information Age quackery
As technology has evolved, particularly with the advent and wide adoption of the internet, it has increasingly become a source of quackery. For example, writing in The New York Times Magazine, Virginia Heffernan criticized WebMD for biasing readers toward drugs that are sold by the site's pharmaceutical sponsors, even when they are unnecessary. She wrote that WebMD "has become permeated with pseudomedicine and subtle misinformation."
| Biology and health sciences | Alternative and traditional medicine | null |
212764 | https://en.wikipedia.org/wiki/Superluminous%20supernova | Superluminous supernova | A super-luminous supernova (SLSN, plural super luminous supernovae or SLSNe) is a type of stellar explosion with a luminosity 10 or more times higher than that of standard supernovae. Like supernovae, SLSNe seem to be produced by several mechanisms, which is readily revealed by their light-curves and spectra. There are multiple models for what conditions may produce an SLSN, including core collapse in particularly massive stars, millisecond magnetars, interaction with circumstellar material (CSM model), or pair-instability supernovae.
The first confirmed superluminous supernova connected to a gamma ray burst was not found until 2003, when GRB 030329 illuminated the Leo constellation. SN 2003dh represented the death of a star 25 times more massive than the Sun, with material being blasted out at over a tenth the speed of light.
Stars with are likely to produce superluminous supernovae.
Classification
Discoveries of many SLSNe in the 21st century showed that not only were they more luminous by an order of magnitude than most supernovae, their remnants were also unlikely to be powered by the typical radioactive decay that is responsible for the observed energies of conventional supernovae.
SLSNe events use a separate classification scheme to distinguish them from the conventional type Ia, type Ib/Ic, and type II supernovae, roughly distinguishing between the spectral signature of hydrogen-rich and hydrogen-poor events.
Hydrogen-rich SLSNe are classified as Type SLSN-II, with observed radiation passing through the changing opacity of a thick expanding hydrogen envelope. Most hydrogen-poor events are classified as Type SLSN-I, with its visible radiation produced from a large expanding envelope of material powered by an unknown mechanism. A third less common group of SLSNe is also hydrogen-poor and abnormally luminous, but clearly powered by radioactivity from 56Ni.
Increasing number of discoveries find that some SLSNe do not fit cleanly into these three classes, so further sub-classes or unique events have been described. Many or all SLSN-I show spectra without hydrogen or helium but have lightcurves comparable to conventional type Ic supernovae, and are now classed as SLSN-Ic. PS1-10afx is an unusually red hydrogen-free SLSN with an extremely rapid rise to a near-record peak luminosity and an unusually rapid decline. PS1-11ap is similar to a type Ic SLSN but has an unusually slow rise and decline.
Astrophysical models
A wide variety of causes have been proposed to explain events that are an order of magnitude or more greater than standard supernovae. The collapsar and CSM (circumstellar material) models are generally accepted and a number of events are well-observed. Other models are still only tentatively accepted or remain entirely theoretical.
Collapsar model
The collapsar model is a type of superluminous supernova that produces a gravitationally collapsed object, or black hole. The word "collapsar", short for "collapsed star", was formerly used to refer to the end product of stellar gravitational collapse, a stellar-mass black hole. The word is now sometimes used to refer to a specific model for the collapse of a fast-rotating star. When core collapse occurs in a star with a core at least around fifteen times the Sun's mass ()—though chemical composition and rotational rate are also significant—the explosion energy is insufficient to expel the outer layers of the star, and it will collapse into a black hole without producing a visible supernova outburst.
A star with a core mass slightly below this level—in the range of —will undergo a supernova explosion, but so much of the ejected mass falls back onto the core remnant that it still collapses into a black hole. If such a star is rotating slowly, then it will produce a faint supernova, but if the star is rotating quickly enough, then the fallback to the black hole will produce relativistic jets. The energy that these jets transfer into the ejected shell renders the visible outburst substantially more luminous than a standard supernova. The jets also beam high energy particles and gamma rays directly outward and thereby produce x-ray or gamma-ray bursts; the jets can last for several seconds or longer and correspond to long-duration gamma-ray bursts, but they do not appear to explain short-duration gamma-ray bursts.
Stars with cores have an approximate total mass of , assuming the star has not undergone significant mass loss. Such a star will still have a hydrogen envelope and will explode as a Type II supernova. Faint Type II supernovae have been observed, but no definite candidates for a Type II SLSN (except type IIn, which are not thought to be jet supernovae). Only the very lowest metallicity population III stars will reach this stage of their life with little mass loss. Other stars, including most of those visible to us, will have had most of their outer layers blown away by their high luminosity and become Wolf-Rayet stars. Some theories propose these will produce either Type Ib or Type Ic supernovae, but none of these events so far has been observed in nature. Many observed SLSNe are likely Type Ic. Those associated with gamma-ray bursts are almost always Type Ic, being very good candidates for having relativistic jets produced by fallback to a black hole. However, not all Type Ic SLSNe correspond to observed gamma-ray bursts but the events would only be visible if one of the jets were aimed towards us.
In recent years, much observational data on long-duration gamma-ray bursts have significantly increased our understanding of these events and made clear that the collapsar model produces explosions that differ only in detail from more or less ordinary supernovae and have energy ranges from approximately normal to around 100 times larger.
A good example of a collapsar SLSN is SN 1998bw, which was associated with the gamma-ray burst GRB 980425. It is classified as a type Ic supernova due to its distinctive spectral properties in the radio spectrum, indicating the presence of relativistic matter.
Circumstellar material model
Almost all observed SLSNe have had spectra similar to either a type Ic or type IIn supernova. The type Ic SLSNe are thought to be produced by jets from fallback to a black hole, but type IIn SLSNe have significantly different light curves and are not associated with gamma-ray bursts. Type IIn supernovae are all embedded in a dense nebula probably expelled from the progenitor star itself, and this circumstellar material (CSM) is thought to be the cause of the extra luminosity. When material expelled in an initial normal supernova explosion meets dense nebular material or dust close to the star, the shockwave converts kinetic energy efficiently into visible radiation. This effect greatly enhances these extended duration and extremely luminous supernovae, even though the initial explosive energy was the same as that of normal supernovae.
Although any supernova type could potentially produce Type IIn SLSNe, theoretical constraints on the surrounding CSM sizes and densities do suggest that it will almost always be produced from the central progenitor star itself immediately prior to the observed supernova event. Such stars are likely candidates of hypergiants or LBVs that appear to be undergoing substantial mass loss, due to Eddington instability, for example, SN2005gl.
Pair-instability supernova
Another type of suspected SLSN is a pair-instability supernova, of which SN 2006gy may possibly be the first observed example. This supernova event was observed in a galaxy about 238 million light years (73 megaparsecs) from Earth.
The theoretical basis for pair-instability collapse has been known for many decades and was suggested as a dominant source of higher mass elements in the early universe as super-massive population III stars exploded. In a pair-instability supernova, the pair production effect causes a sudden pressure drop in the star's core, leading to a rapid partial collapse. Gravitational potential energy from the collapse causes runaway fusion of the core which entirely disrupts the star, leaving no remnant.
Models show that this phenomenon only happens in stars with extremely low metallicity and masses between about 130 and 260 times the Sun, making them extremely unlikely in the local universe. Although originally expected to produce SLSN explosions hundreds of times greater than a normal supernova, current models predict that they actually produce luminosities ranging from about the same as a normal core collapse supernova to perhaps 50 times brighter, although remaining bright for much longer.
Magnetar energy release
Models of the creation and subsequent spin-down of a magnetar yield much higher luminosities than regular supernova events and match the observed properties of at least some SLSNe. In cases where pair-instability supernova may not be a good fit for explaining a SLSN, a magnetar explanation is more plausible.
Other models
There are still models for SLSN explosions produced from binary systems, white dwarf or neutron stars in unusual arrangements or undergoing mergers, and some of these are proposed to account for some observed gamma-ray bursts.
| Physical sciences | Stellar astronomy | Astronomy |
212786 | https://en.wikipedia.org/wiki/Maltose | Maltose | Maltose ( or ), also known as maltobiose or malt sugar, is a disaccharide formed from two units of glucose joined with an α(1→4) bond. In the isomer isomaltose, the two glucose molecules are joined with an α(1→6) bond. Maltose is the two-unit member of the amylose homologous series, the key structural motif of starch. When beta-amylase breaks down starch, it removes two glucose units at a time, producing maltose. An example of this reaction is found in germinating seeds, which is why it was named after malt. Unlike sucrose, it is a reducing sugar.
History
Maltose was discovered by Augustin-Pierre Dubrunfaut, although this discovery was not widely accepted until it was confirmed in 1872 by Irish chemist and brewer Cornelius O'Sullivan. Its name comes from malt, combined with the suffix '-ose' which is used in names of sugars.
Structure and nomenclature
Carbohydrates are generally divided into monosaccharides, oligosaccharides, and polysaccharides depending on the number of sugar subunits. Maltose, with two sugar units, is a disaccharide, which falls under oligosaccharides. Glucose is a hexose: a monosaccharide containing six carbon atoms. The two glucose units are in the pyranose form and are joined by an O-glycosidic bond, with the first carbon (C1) of the first glucose linked to the fourth carbon (C4) of the second glucose, indicated as (1→4). The link is characterized as α because the glycosidic bond to the anomeric carbon (C1) is in the opposite plane from the substituent in the same ring (C6 of the first glucose). If the glycosidic bond to the anomeric carbon (C1) were in the same plane as the substituent, it would be classified as a β(1→4) bond, and the resulting molecule would be cellobiose. The anomeric carbon (C1) of the second glucose molecule, which is not involved in a glycosidic bond, could be either an α- or β-anomer depending on the bond direction of the attached hydroxyl group relative to the substituent of the same ring, resulting in either α-maltose or β-maltose.
An isomer of maltose is isomaltose. This is similar to maltose but instead of a bond in the α(1→4) position, it is in the α(1→6) position, the same bond that is found at the branch points of glycogen and amylopectin.
Properties
Like glucose, maltose is a reducing sugar, because the ring of one of the two glucose units can open to present a free aldehyde group; the other one cannot because of the nature of the glycosidic bond. Maltose can be broken down to glucose by the maltase enzyme, which catalyses the hydrolysis of the glycosidic bond.
Maltose in aqueous solution exhibits mutarotation, because the α and β isomers that are formed by the different conformations of the anomeric carbon have different specific rotations, and in aqueous solutions, these two forms are in equilibrium. Maltose can easily be detected by the Woehlk test or Fearon's test on methylamine.
It has a sweet taste, but is only about 30–60% as sweet as sugar, depending on the concentration. A 10% solution of maltose is 35% as sweet as sucrose.
Sources and absorption
Maltose is a malt component, a substance obtained when the grain is softened in water and germinates. It is also present in highly variable quantities in partially hydrolyzed starch products like maltodextrin, corn syrup and acid-thinned starch.
Outside of plants, maltose is also (likely) found in sugarbag.
In humans, maltose is broken down by various maltase enzymes, providing two glucose molecules that can be further processed: either broken down to provide energy, or stored as glycogen. The lack of the sucrase-isomaltase enzyme in humans causes sucrose intolerance, but complete maltose intolerance is extremely rare because there are four different maltase enzymes.
| Biology and health sciences | Carbohydrates | Biology |
212855 | https://en.wikipedia.org/wiki/Memory%20card | Memory card | A memory card is an electronic data storage device used for storing digital information, typically using flash memory. These are commonly used in digital portable electronic devices, such as digital cameras as well as in many early games consoles such as the Nintendo Wii. They allow adding memory to such devices using a card in a socket instead of protruding USB flash drives.
Common types of flash memory card include SD cards (including microSD), Sony's Memory Stick and CompactFlash. , SD cards are the most common type of memory cards.
History
The basis for memory card technology is flash memory. It was invented by Fujio Masuoka at Toshiba in 1980 and commercialized by Toshiba in 1987.
The development of memory cards was driven in the 1980s by the need for an alternative to floppy disk drives that had lower power consumption, had less weight and occupied less volume in laptops. Some were also marketed as a lower cost alternative to ROM cartridges. Several competing and incompatible memory card formats were developed by several vendors, such as for example the Bee Card, Astron SoftCards, Sega Cards, NEC UltraLite memory cards, and the Mitsubishi Melcard which came in variants using 60 and 50 connector pins. The Sega Card was developed as a cheaper alternative to game cartridges. Some memory cards were used for memory expansion in laptops.
JEIDA, the Japan Electronic Industry Development Association, began to work on a standard for memory cards in 1985, and developed the JEIDA memory card in 1986.
The Personal Computer Memory Card International Association (PCMCIA) was an industry association created in 1989 to promote a standard for memory cards in PCs, and worked closely with JEIDA, adopting their 68 pin connector design. The specification for PCMCIA type I cards, later renamed PC Cards, was first released in 1990, and unified the JEIDA memory card standard with the PC Card standard. This format later included support for other devices besides memory cards. PC Card was among the first commercial memory card formats to come out, but is mainly used in industrial applications and to connect I/O devices such as modems.
Some early memory cards used SRAM as a storage medium, which required a lithium battery to keep the contents in the SRAM. These cards were faster than their flash counterparts. Some of the first PCMCIA cards had capacities of 1 to 5 MB and cost US$100 per MB. Other early cards such as the Bee Card contained non-modifiable ROM, Write once read many EPROM or rewriteable EEPROM memory. In 1992, SanDisk introduced FlashDisk, a PCMCIA card and one of the first memory cards that did not require battery power to retain its contents, as it used flash memory.
In 1994, memory card formats smaller than the PC Card arrived. The first one was CompactFlash and later SmartMedia and Miniature Card. The desire for smaller cards for cell-phones, PDAs, and compact digital cameras drove a trend that left the previous generation of "compact" cards looking big. In 2000 the SD card was announced. SD was envisioned as a single memory card format for several kinds of electronic devices, that could also function as an expansion slot for adding new capabilities for a device. In 2001, SmartMedia alone captured 50% of the digital camera market and CF had captured the professional digital camera market.
However, by 2005, SD and similar MMC cards had nearly taken over SmartMedia's spot, though not to the same level and with stiff competition coming from Memory Stick variants, as well as CompactFlash. In industrial and embedded fields, even the venerable PC card (PCMCIA) memory cards still manage to maintain a niche, while in mobile phones and PDAs, the memory card has become smaller.
Initially memory cards were expensive, costing US$3 per megabyte of capacity in 2001; this led to the development of miniaturized rotating disk memory devices such as the Microdrive, PocketZip and Dataplay. The Microdrive had higher capacities than memory cards at the time. All three concepts became obsolete once flash memory prices became lower and their capacities became higher by 2006.
New products of Sony (previously only using Memory Stick) and Olympus (previously only using XD-Card) have been offered with an additional SD-Card slot beginning in 2010. Effectively the format war has turned in SD-Card's favor.
Data table of selected memory card formats
Overview of all memory card types
PCMCIA ATA Type I Card (PC Card ATA Type I)
PCMCIA Type II, Type III cards
CompactFlash Card (Type I), CompactFlash High-Speed
CompactFlash Type II, CF+(CF2.0), CF3.0
Microdrive
CFexpress
MiniCard (Miniature Card) (max 64 MB / 64 MiB)
SmartMedia Card (SSFDC) (max 128 MB) (3.3 V,5 V)
xD-Picture Card, xD-Picture Card Type M
Memory Stick, MagicGate Memory Stick (max 128 MB); Memory Stick Select, MagicGate Memory Stick Select ("Select" means: 2x128 MB with A/B switch)
SecureMMC
Secure Digital (SD Card), Secure Digital High-Speed, Secure Digital Plus/Xtra/etc (SD with USB connector)
miniSD card
microSD card (aka Transflash, T-Flash, TF)
SDHC
WiFi SD Cards (SD Card With WiFi Card Built in) Powered by Device. (Eye-Fi, WiFi SD, Flash Air)
Nano Memory (NM) card
MU-Flash (Mu-Card) (Mu-Card Alliance of OMIA)
C-Flash
SIM card (Subscriber Identity Module)
Smart card (ISO/IEC 7810, ISO/IEC 7816 card standards, etc.)
UFC (USB FlashCard) (uses USB)
FISH Universal Transportable Memory Card Standard (uses USB)
Intelligent Stick (iStick, a USB-based flash memory card with MMS)
SxS (S-by-S) memory card, a new memory card specification developed by Sandisk and Sony. SxS complies to the ExpressCard industry standard.
Nexflash Winbond Serial Flash Module (SFM) cards, size range 1 MB, 2 MB and 4 MB.
Comparison
Video game consoles
Many older video game consoles used memory cards to hold saved game data. Cartridge-based systems primarily used battery-backed volatile RAM within each individual cartridge to hold saves for that game. Cartridges without this RAM may have used a password system, or would not save progress at all. The Neo Geo AES, released in 1990 by SNK, was the first video game console able to use a memory card. AES memory cards were also compatible with Neo Geo MVS arcade cabinets, allowing players to migrate saves between home and arcade systems and vice versa. Memory cards became commonplace when home consoles moved to read-only optical discs for storing the game program, beginning with systems such as the TurboGrafx-CD and Sega-CD.
Until the sixth generation of video game consoles, memory cards were based on proprietary formats; Later systems used established industry formats for memory cards, such as FAT32.
Home consoles commonly use hard disk drive storage for saved games and allow the use of USB flash drives or other card formats via a memory card reader to transport game saves and other game information. Though some consoles have implemented cloud storage saving, most portable gaming systems still rely on custom memory cartridges to store program data, due to their low power consumption, smaller physical size and reduced mechanical complexity.
| Technology | Data storage | null |
212980 | https://en.wikipedia.org/wiki/Order%20of%20operations | Order of operations | In mathematics and computer programming, the order of operations is a collection of rules that reflect conventions about which operations to perform first in order to evaluate a given mathematical expression.
These rules are formalized with a ranking of the operations. The rank of an operation is called its precedence, and an operation with a higher precedence is performed before operations with lower precedence. Calculators generally perform operations with the same precedence from left to right, but some programming languages and calculators adopt different conventions.
For example, multiplication is granted a higher precedence than addition, and it has been this way since the introduction of modern algebraic notation. Thus, in the expression , the multiplication is performed before addition, and the expression has the value , and not . When exponents were introduced in the 16th and 17th centuries, they were given precedence over both addition and multiplication and placed as a superscript to the right of their base. Thus and .
These conventions exist to avoid notational ambiguity while allowing notation to remain brief. Where it is desired to override the precedence conventions, or even simply to emphasize them, parentheses ( ) can be used. For example, forces addition to precede multiplication, while forces addition to precede exponentiation. If multiple pairs of parentheses are required in a mathematical expression (such as in the case of nested parentheses), the parentheses may be replaced by other types of brackets to avoid confusion, as in .
These rules are meaningful only when the usual notation (called infix notation) is used. When functional or Polish notation are used for all operations, the order of operations results from the notation itself.
Conventional order
The order of operations, that is, the order in which the operations in an expression are usually performed, results from a convention adopted throughout mathematics, science, technology and many computer programming languages. It is summarized as:
Parentheses
Exponentiation
Multiplication and division
Addition and subtraction
This means that to evaluate an expression, one first evaluates any sub-expression inside parentheses, working inside to outside if there is more than one set. Whether inside parenthesis or not, the operation that is higher in the above list should be applied first. Operations of the same precedence are conventionally evaluated from left to right.
If each division is replaced with multiplication by the reciprocal (multiplicative inverse) then the associative and commutative laws of multiplication allow the factors in each term to be multiplied together in any order. Sometimes multiplication and division are given equal precedence, or sometimes multiplication is given higher precedence than division; see below. If each subtraction is replaced with addition of the opposite (additive inverse), then the associative and commutative laws of addition allow terms to be added in any order.
The radical symbol is traditionally extended by a bar (called vinculum) over the radicand (this avoids the need for parentheses around the radicand). Other functions use parentheses around the input to avoid ambiguity. The parentheses can be omitted if the input is a single numerical variable or constant, as in the case of and . Traditionally this convention extends to monomials; thus, and even , but , because is not a monomial. However, this convention is not universally understood, and some authors prefer explicit parentheses. Some calculators and programming languages require parentheses around function inputs, some do not.
Symbols of grouping can be used to override the usual order of operations. Grouped symbols can be treated as a single expression. Symbols of grouping can be removed using the associative and distributive laws, also they can be removed if the expression inside the symbol of grouping is sufficiently simplified so no ambiguity results from their removal.
Examples
Multiplication before addition:
Parenthetical subexpressions are evaluated first:
Exponentiation before multiplication, multiplication before subtraction:
When an expression is written as a superscript, the superscript is considered to be grouped by its position above its base:
The operand of a root symbol is determined by the overbar:
A horizontal fractional line also acts as a symbol of grouping:
Parentheses can be nested, and should be evaluated from the inside outward. For legibility, outer parentheses can be made larger than inner parentheses. Alternately, other grouping symbols, such as curly braces } or square brackets , are sometimes used along with parentheses . For example:
Special cases
Unary minus sign
There are differing conventions concerning the unary operation (usually pronounced "minus"). In written or printed mathematics, the expression −32 is interpreted to mean .
In some applications and programming languages, notably Microsoft Excel, PlanMaker (and other spreadsheet applications) and the programming language bc, unary operations have a higher priority than binary operations, that is, the unary minus has higher precedence than exponentiation, so in those languages −32 will be interpreted as . This does not apply to the binary minus for example in Microsoft Excel while the formulas =-2^2, =-(2)^2 and =0+-2^2 return 4, the formulas =0-2^2 and =-(2^2) return −4.
Mixed division and multiplication
There is no universal convention for interpreting an expression containing both division denoted by '÷' and multiplication denoted by '×'. Proposed conventions include assigning the operations equal precedence and evaluating them from left to right, or equivalently treating division as multiplication by the reciprocal and then evaluating in any order; evaluating all multiplications first followed by divisions from left to right; or eschewing such expressions and instead always disambiguating them by explicit parentheses.
Beyond primary education, the symbol '÷' for division is seldom used, but is replaced by the use of algebraic fractions, typically written vertically with the numerator stacked above the denominator – which makes grouping explicit and unambiguous – but sometimes written inline using the slash or solidus symbol, '/'.
Multiplication denoted by juxtaposition (also known as implied multiplication) creates a visual unit and has higher precedence than most other operations. In academic literature, when inline fractions are combined with implied multiplication without explicit parentheses, the multiplication is conventionally interpreted as having higher precedence than division, so that e.g. is interpreted to mean rather than . For instance, the manuscript submission instructions for the Physical Review journals directly state that multiplication has precedence over division, and this is also the convention observed in physics textbooks such as the Course of Theoretical Physics by Landau and Lifshitz and mathematics textbooks such as Concrete Mathematics by Graham, Knuth, and Patashnik. However, some authors recommend against expressions such as , preferring the explicit use of parenthesis .
More complicated cases are more ambiguous. For instance, the notation could plausibly mean either
or . Sometimes interpretation depends on context. The Physical Review submission instructions recommend against expressions of the form ; more explicit expressions or are unambiguous.
This ambiguity has been the subject of Internet memes such as "", for which there are two conflicting interpretations: 8 ÷ [2 · (2 + 2)] = 1 and (8 ÷ 2) · (2 + 2) = 16. Mathematics education researcher Hung-Hsi Wu points out that "one never gets a computation of this type in real life", and calls such contrived examples "a kind of Gotcha! parlor game designed to trap an unsuspecting person by phrasing it in terms of a set of unreasonably convoluted rules."
Serial exponentiation
If exponentiation is indicated by stacked symbols using superscript notation, the usual rule is to work from the top down:
which typically is not equal to (ab)c. This convention is useful because there is a property of exponentiation that (ab)c = abc, so it's unnecessary to use serial exponentiation for this.
However, when exponentiation is represented by an explicit symbol such as a caret (^) or arrow (↑), there is no common standard. For example, Microsoft Excel and computation programming language MATLAB evaluate a^b^c as (ab)c, but Google Search and Wolfram Alpha as a(bc). Thus 4^3^2 is evaluated to 4,096 in the first case and to 262,144 in the second case.
Mnemonics
Mnemonic acronyms are often taught in primary schools to help students remember the order of operations. The acronym PEMDAS, which stands for Parentheses, Exponents, Multiplication/Division, Addition/Subtraction, is common in the United States and France. Sometimes the letters are expanded into words of a mnemonic sentence such as "Please Excuse My Dear Aunt Sally". The United Kingdom and other Commonwealth countries may use BODMAS (or sometimes BOMDAS), standing for Brackets, Of, Division/Multiplication, Addition/Subtraction, with "of" meaning fraction multiplication. Sometimes the O is instead expanded as Order, meaning exponent or root, or replaced by I for Indices in the alternative mnemonic BIDMAS. In Canada and New Zealand BEDMAS is common.
In Germany, the convention is simply taught as , "dot operations before line operations" referring to the graphical shapes of the taught operator signs , , and , .
These mnemonics may be misleading when written this way. For example, misinterpreting any of the above rules to mean "addition first, subtraction afterward" would incorrectly evaluate the expression as , while the correct evaluation is . These values are different when .
Mnemonic acronyms have been criticized for not developing a conceptual understanding of the order of operations, and not addressing student questions about its purpose or flexibility. Students learning the order of operations via mnemonic acronyms routinely make mistakes, as do some pre-service teachers. Even when students correctly learn the acronym, a disproportionate focus on memorization of trivia crowds out substantive mathematical content. The acronym's procedural application does not match experts' intuitive understanding of mathematical notation: mathematical notation indicates groupings in ways other than parentheses or brackets and a mathematical expression is a tree-like hierarchy rather than a linearly "ordered" structure; furthermore, there is no single order by which mathematical expressions must be simplified or evaluated and no universal canonical simplification for any particular expression, and experts fluently apply valid transformations and substitutions in whatever order is convenient, so learning a rigid procedure can lead students to a misleading and limiting understanding of mathematical notation.
Calculators
Different calculators follow different orders of operations. Many simple calculators without a stack implement chain input, working in button-press order without any priority given to different operations, give a different result from that given by more sophisticated calculators. For example, on a simple calculator, typing 1 + 2 × 3 = yields 9, while a more sophisticated calculator will use a more standard priority, so typing 1 + 2 × 3 = yields 7.
Calculators may associate exponents to the left or to the right. For example, the expression a^b^c is interpreted as a(bc) on the TI-92 and the TI-30XS MultiView in "Mathprint mode", whereas it is interpreted as (ab)c on the TI-30XII and the TI-30XS MultiView in "Classic mode".
An expression like 1/2x is interpreted as 1/(2x) by TI-82, as well as many modern Casio calculators (configurable on some like the fx-9750GIII), but as (1/2)x by TI-83 and every other TI calculator released since 1996, as well as by all Hewlett-Packard calculators with algebraic notation. While the first interpretation may be expected by some users due to the nature of implied multiplication, the latter is more in line with the rule that multiplication and division are of equal precedence.
When the user is unsure how a calculator will interpret an expression, parentheses can be used to remove the ambiguity.
Order of operations arose due to the adaptation of infix notation in standard mathematical notation, which can be notationally ambiguous without such conventions, as opposed to postfix notation or prefix notation, which do not need orders of operations. Hence, calculators utilizing Reverse Polish notation (RPN) using a stack to enter expressions in the correct order of precedence do not need parentheses or any possibly model-specific order of execution.
Programming languages
Most programming languages use precedence levels that conform to the order commonly used in mathematics, though others, such as APL, Smalltalk, Occam and Mary, have no operator precedence rules (in APL, evaluation is strictly right to left; in Smalltalk, it is strictly left to right).
Furthermore, because many operators are not associative, the order within any single level is usually defined by grouping left to right so that 16/4/4 is interpreted as rather than ; such operators are referred to as "left associative". Exceptions exist; for example, languages with operators corresponding to the cons operation on lists usually make them group right to left ("right associative"), e.g. in Haskell, 1:2:3:4:[] == 1:(2:(3:(4:[]))) == [1,2,3,4].
Dennis Ritchie, creator of the C language, said of the precedence in C (shared by programming languages that borrow those rules from C, for example, C++, Perl and PHP) that it would have been preferable to move the bitwise operators above the comparison operators. Many programmers have become accustomed to this order, but more recent popular languages like Python and Ruby do have this order reversed. The relative precedence levels of operators found in many C-style languages are as follows:
Examples:
!A + !B is interpreted as (!A) + (!B)
++A + !B is interpreted as (++A) + (!B)
A + B * C is interpreted as A + (B * C)
A || B && C is interpreted as A || (B && C)
A && B == C is interpreted as A && (B == C)
A & B == C is interpreted as A & (B == C)
(In Python, Ruby, PARI/GP and other popular languages, A & B == C is interpreted as (A & B) == C.)
Source-to-source compilers that compile to multiple languages need to explicitly deal with the issue of different order of operations across languages. Haxe for example standardizes the order and enforces it by inserting brackets where it is appropriate.
The accuracy of software developer knowledge about binary operator precedence has been found to closely follow their frequency of occurrence in source code.
History
The Order of Operations emerged progressively over centuries. The rule that multiplication has precedence over addition was incorporated into the development of algebraic notation in the 1600s, since the distributive property implies this as a natural hierarchy. As recently as the 1920s, the historian of mathematics Florian Cajori identifies disagreement about whether multiplication should have precedence over division, or whether they should be treated equally. The term "order of operations" and the "PEMDAS/BEDMAS" mnemonics were formalized only in the late 19th or early 20th century, as demand for standardized textbooks grew.
Ambiguity about issues such as whether implicit multiplication takes precedence over explicit multiplication and division in such expressions as a/2b, which could be interpreted as a/(2b) or (a/2)*b, imply that the conventions are not yet completely stable.
| Mathematics | Basics | null |
212989 | https://en.wikipedia.org/wiki/Western%20jackdaw | Western jackdaw | The western jackdaw (Coloeus monedula), also known as the Eurasian jackdaw, the European jackdaw, or simply the jackdaw, is a passerine bird in the crow family. Found across Europe, western Asia and western North Africa; it is mostly resident, although northern and eastern populations migrate south in the winter. Four subspecies are recognised, which differ mainly in the colouration of the plumage on the head and nape. Linnaeus first described it formally, giving it the name Corvus monedula. The common name derives from the word jack, denoting "small", and daw, a less common synonym for "jackdaw", and the native English name for the bird.
Measuring in length, the western jackdaw is a black-plumaged bird with a grey nape and distinctive pale-grey irises. It is gregarious and vocal, living in small groups with a complex social structure in farmland, open woodland, on coastal cliffs, and in urban settings. Like its relatives, jackdaws are intelligent birds, and have been observed using tools. An omnivorous and opportunistic feeder, it eats a wide variety of plant material and invertebrates, as well as food waste from urban areas. Western jackdaws are monogamous and build simple nests of sticks in cavities in trees, cliffs, or buildings. About five pale blue or blue-green eggs with brown speckles are laid and incubated by the female. The young fledge in four to five weeks.
Systematics
Etymology
The western jackdaw was one of the many species originally described by Carl Linnaeus in his 18th century work Systema Naturae. Owing to its supposed fondness for picking up coins, Linnaeus gave it the binomial name Corvus monedula, choosing the specific name mǒnēdŭla, which is derived from moneta, the Latin stem of the word "money". Jackdaws are sometimes placed in the genus Coloeus, from the Ancient Greek κολοιός (koloios) for jackdaw, though most subsequent works have retained the two jackdaw species in Corvus.
The original Old English words ċēo and ċeahhe (pronounced with initial ch) gave modern English "chough"; Chaucer sometimes used this word to refer to the western jackdaw, as did Shakespeare in Hamlet although there has been debate about which species he was referring to. This onomatopoeic name, based on the western jackdaw's call, now refers to corvids of the genus Pyrrhocorax; the red-billed chough (Pyrrhocorax pyrrhocorax), formerly particularly common in Cornwall, became known initially as the "Cornish chough" and then just the "chough", the name transferring from one species to the other.
The common name jackdaw first appeared in the 16th century, and is thought to be a compound of the forename Jack, used in animal names to signify a small form (e.g. jack snipe), and the archaic native English word daw. Formerly, western jackdaws were simply called "daws". The metallic chyak call may be the origin of the jack part of the common name, but this is not supported by the Oxford English Dictionary. Daw, first used for the bird in the 15th century, is held by the Oxford English Dictionary to be derived from the postulated Old English dawe, citing the cognates in Old High German tāha, Middle High German tāhe or tāchele, and modern German Dahle or Dohle, and dialectal Tach, Dähi, Däche and Dacha.
Names in English dialects are numerous. Scottish and north English dialects have included ka or kae since the 14th century. The Midlands form of this word was co or coo. Caddow is potentially a compound of ka and dow, a variant of daw. Other dialectal or obsolete names include caddesse, cawdaw, caddy, chauk, college-bird, jackerdaw, jacko, ka-wattie, chimney-sweep bird (from their nesting propensities), and sea-crow (from the frequency with which they are found on coasts). It was also frequently known quasi-nominally as Jack.
An archaic collective noun for a group of jackdaws is a "clattering". Another name for a flock is a "train".
Taxonomy
A study in 2000 found that the genetic distance between western jackdaws and the other members of Corvus was greater than that within the rest of the genus. This led Pamela Rasmussen to reinstate the genus name Coloeus, created by Johann Kaup in 1829, in her Birds of South Asia (2005), a treatment also used in a 1982 systematic list in German by Hans Edmund Wolters. A study of corvid phylogeny undertaken in 2007 compared DNA sequences in the mitochondrial control region of several corvids. It found that the western jackdaw, and the closely related Daurian jackdaw (C. dauuricus) of eastern Russia and China, were basal to the core Corvus clade. The names Coloeus monedula and Coloeus dauuricus have since been adopted by the International Ornithological Congress in their official list. The two species of jackdaw have been reported to hybridise in the Altai Mountains, southern Siberia, and Mongolia. Analysis of the mitochondrial DNA of specimens of the two species from their core ranges show them to be genetically distinct.
Subspecies
There are four recognized subspecies of the western jackdaw. All European subspecies intergrade where their populations meet. C. m. monedula intergrades into C. m. soemmerringii in a transition zone running from Finland south across the Baltic and eastern Poland to Romania and Croatia.
The Nordic jackdaw (C. m. monedula) (Linnaeus, 1758), the nominate subspecies, is found in eastern Europe. Its range extends across Scandinavia, from southern Finland south to Esbjerg and Haderslev in Denmark, through eastern Germany and Poland, and south across eastern central Europe to the Carpathian Mountains and northwestern Romania, Vojvodina in northern Serbia, and Slovenia. It breeds in south-eastern Norway, southern Sweden, and northern and eastern Denmark, with occasional wintering in England and France. It has been recorded as a rare vagrant to Spain. It has a pale nape and sides of the neck, a dark throat, and a light grey partial collar of variable extent.
The Western Eurasian jackdaw (C. m. spermologus) (Vieillot, 1817) occurs in western, central and southern Europe and North Africa, from the British Isles, the Netherlands and the Rhineland in the north, through western Switzerland into Italy in the southeast, and the Iberian peninsula and Morocco in the south. It winters in the Canary Islands and Corsica. The name "spermologus" comes from the Greek σπερμολόγος, a picker of seeds. It is darker in colour than the other subspecies and lacks the whitish border at the base of the grey nape.
The Eastern Eurasian jackdaw (C. m. soemmerringii) (Fischer von Waldheim, 1811) is found in northeastern Europe and northern and central Asia from the former Soviet Union to Lake Baikal and northwestern Mongolia and south to Turkey, Israel and the eastern Himalayas. Its southwestern limits are Serbia and southern Romania. It winters in Iran and northern India (Kashmir). Johann Fischer von Waldheim described this taxon as Corvus soemmerringii in 1811, noting its differences from populations in western Europe. Its subspecific name was given in honour of the German anatomist Samuel Thomas von Sömmerring. It is distinguished by the nape and the sides of the neck being paler, creating a contrasting black crown and lighter grey part collar.
The Algerian jackdaw (C. m. cirtensis) (Rothschild and Hartert, 1912) is found in Morocco and Algeria in Northwest Africa. It was also formerly found in Tunisia. The name "cirtensis" refers to the ancient city of Cirta in Numidia. The plumage is duller and more uniformly dark grey than the other subspecies, with the paler nape less distinct.
Description
The western jackdaw measures in length and weighs around . Most of the plumage is a shiny black, with a purple (in subspecies monedula and spermologus) or blue (in subspecies cirtensis and soemmerringii) sheen on the crown, forehead, and secondaries, and a green-blue sheen on the throat, primaries, and tail. The cheeks, nape and neck are light grey to greyish-silver, and the underparts are slate-grey. The legs are black, as is the short stout bill, the length of which is about 75% of the length of the rest of the head. There are rictal bristles covering around 40% of the maxilla and 25% of the lower mandible. The irises of adults are greyish or silvery white while those of juveniles are light blue, becoming brownish before whitening at around one year of age. The sexes look alike, though the head and neck plumage of male birds fades more with age and wear, particularly just before moulting. Western jackdaws undergo a complete moult from June to September in the western parts of their range, and a month later in the east. The purplish sheen of the cap is most prominent just after moulting.
Immature birds have duller and less demarcated plumage. The head is a sooty black, sometimes with a faint greenish sheen and brown feather bases visible; the back and side of the neck are dark grey and the underparts greyish or sooty black. The tail has narrower feathers and a greenish sheen.
There is very little geographic variation in size. The main differences are the presence or absence of a whitish partial collar at the base of the nape, the variations in the shade of the nape and the tone of the underparts. Populations in central Asia have slightly larger wings and western populations have a slightly heavier bill. Body colour becomes darker further north, in mountain regions and humid climates, and paler elsewhere. However, individual variation, particularly in juveniles and also during the months before moulting, can often be greater than geographic differences.
A skilled flyer, the western jackdaw can manoeuvre tightly as well as tumble and glide. It has characteristic jerky wing beats when flying, though these are not evident when birds are migrating. Wind tunnel experiments show that the preferred gliding speed is between per second and that the wingspan decreases as the bird flies faster. On the ground, western jackdaws have an upright posture and strut briskly, their short legs giving them a rapid gait. They feed with their heads held down or horizontally.
Within its range, the western jackdaw is unmistakable; its short bill and grey nape are distinguishing features. From a distance, it can be confused with a rook (Corvus frugilegus), or when in flight, with a pigeon or chough. Flying western jackdaws are distinguishable from other corvids by their smaller size, faster and deeper wingbeats and proportionately narrower and less fingered wing tips. They also have shorter, thicker necks, much shorter bills and frequently fly in tighter flocks. They can be distinguished from choughs by their uniformly grey underwings and their black beaks and legs. The western jackdaw is very similar in morphology, behaviour, and calls to the Daurian jackdaw, with which its range overlaps in western Asia. Adults are readily distinguished, since the Daurian has a pied plumage, but immature birds are much more similar, both species having dark plumage and dark eyes. The Daurian tends to be darker, with a less contrasting nape than the Western.
Vocalisations
Western jackdaws are voluble birds. The main call, frequently given in flight, is a metallic and squeaky chyak-chyak or kak-kak. This is a contact or greeting call. A feeding call made by adults to call young, or males when offering food to their mates, has been transcribed as kiaw or kyow. Females in return give a more drawn out version when begging for food from males, written as kyaay, tchaayk or giaaaa. Perched birds often chatter together, and before settling for the night, large roosting flocks make a cackling noise. Western jackdaws also have a hoarse, drawn-out alarm call, arrrrr or kaaaarr, used when warning of predators or when mobbing them. Nestlings begin making a soft cheep at about a week of age. As they grow, their voice becomes louder until their call is a penetrating screech around day 18. After this, the voice deepens and softens. From day 25, the young cease calling and become silent if they hear an unfamiliar noise. The European jackdaw can be trained to imitate human speech.
Distribution and habitat
The western jackdaw is found from Northwest Africa through all of Europe, except for the subarctic north, and eastwards through central Asia to the eastern Himalayas and Lake Baikal. To the east, it occurs throughout Turkey, the Caucasus, Iran, Iraq, Afghanistan, Pakistan, and northwestern India. However, it is regionally extinct in Malta and Tunisia. The range is vast, with an estimated global extent between . It has a large global population, with an estimated 15.6 to 45 million individuals in Europe alone. Censuses of bird populations in marginal uplands in Great Britain show that western jackdaws greatly increased in numbers between the 1970s and 2010, although this increase may be related to recovery from previous periods when they were regarded as pests. The UK population was estimated at 2.5 million individuals in 1998, up from 780,000 in 1970.
Most populations are resident, but the northern and eastern populations are more migratory, relocating to wintering areas between September and November and returning between February and early May. Their range expands northwards into Russia to Siberia during summer and retracts in winter. They are vagrants to the Faroe Islands, particularly in the winter and spring, and occasionally to Iceland. Elsewhere, western jackdaws congregate over winter in the Ural Valley in northwestern Kazakhstan, the northern Caspian, and the Tian Shan region of western China. They are winter visitors to the Quetta Valley in western Pakistan, and are winter vagrants to Lebanon, where they were first recorded in 1962. In Syria, they are winter vagrants and rare residents with some confirmed breeding taking place. The subspecies soemmerringii occurs in south-central Siberia and extreme northwestern China and is accidental to Hokkaido, Japan. A small number of western jackdaws reached northeastern North America in the 1980s and have been found from Atlantic Canada to Pennsylvania. They have also occurred as vagrants in Gibraltar, Mauritania, and Saint Pierre and Miquelon, and one is reported to have been seen in Egypt.
Western jackdaws inhabit wooded steppes, pastures, cultivated land, coastal cliffs, and towns. They thrive when forested areas are cleared and converted to fields and open areas. Habitats with a mix of large trees, buildings, and open ground are preferred; open fields are left to the rook, and more wooded areas to the Eurasian jay (Garrulus glandarius). Along with other corvids such as the rook, common raven (Corvus corax), and hooded crow (C. cornix), some western jackdaws spend the winter in urban parks; populations measured in three urban parks in Warsaw show increases from October to December, possibly due to western jackdaws migrating there from areas further north. The same data from Warsaw, collected from 1977 to 2003, showed that the wintering western jackdaw population had increased four-fold. The cause of the increase is unknown, but a reduction in the number of rooks may have benefited the species locally, or rooks overwintering in Belarus may have caused western jackdaws to relocate to Warsaw.
Behaviour
Generally wary of people in the forest or countryside, western jackdaws are much tamer in urban areas.
Highly gregarious, western jackdaws are generally seen in flocks of varying sizes, though males and females pair-bond for life and pairs stay together within flocks. Flocks increase in size in autumn and birds congregate at dusk for communal roosting, with up to several thousand individuals gathering at one site. At Uppsala, Sweden, 40,000 birds have been recorded at a single winter roost with mated pairs often settling together for the night. Western jackdaws frequently congregate with hooded crows or rooks, the latter particularly when migrating or roosting. They have been recorded foraging with the common starling (Sturnus vulgaris), Northern lapwing (Vanellus vanellus), and common gull (Larus canus) in northwestern England. Flocks are targets of coordinated hunting by pairs of lanner falcons (Falco biarmicus), although larger groups are more able to elude the predators. Western jackdaws sometimes mob and drive off larger birds such as European magpies, common ravens, or Egyptian vultures (Neophron percnopterus); one gives an alarm call which alerts its conspecifics to gather and attack as a group. Occasionally, a sick or injured western jackdaw is mobbed until it is killed.
In his book King Solomon's Ring, Konrad Lorenz described and analysed the complex social interactions in a western jackdaw flock that lived around his house in Altenberg, Austria. He ringed them for identification and caged them in the winter to prevent their annual migration. He found that the birds have a linear hierarchical group structure, with higher-ranked individuals dominating lower-ranked birds, and pair-bonded birds sharing the same rank. Young males establish their individual status before pairing with females. Upon pairing, the female assumes the same social position as her partner. Unmated females are the lowest members in the pecking order, and are the last to have access to food and shelter. Lorenz noted one case in which a male, absent during the dominance struggles and pair bondings, returned to the flock, became the dominant male, and chose one of two unpaired females for a mate. This female immediately assumed a dominant position in the social hierarchy and demonstrated this by pecking others. According to Lorenz, the most significant factor in social behaviour was the immediate and intuitive grasp of the new hierarchy by each of the western jackdaws in the flock.
Social displays
Social hierarchy in western jackdaw flocks is determined by supplanting, fighting, and threat displays—several of which have been described. In the bill-up posture, the western jackdaw tilts its bill and head upwards and sleeks its plumage. Indicating both appeasement and assertiveness, the posture is used by birds intending to enter feeding flocks. A bill-down posture is another commonly used agonistic behaviour. In this display, a bird lowers its bill and erects its nape and head feathers, and sometimes slightly lifts its wings. Western jackdaws often face off in this posture until one backs down or a fight ensues. In the forward-threat posture, a bird holds its body horizontally and thrusts its head forwards. In intense versions, the bird ruffles its feathers and spreads or raises its tail and wings. This extreme is seen when facing off over nests or females. In the defensive-threat posture, the bird lowers its head and bill, spreads its tail and ruffles its feathers. Supplanting is where one bird moves in and displaces another from a perch-site. The second bird usually retreats without resorting to a fight. Western jackdaws fight by launching themselves at each other feet-first and then wrestling with their feet intertwined and pecking at each other. Other individuals gather and call noisily.
Western jackdaws entreat their partners to preen them by showing their nape and ruffling their head feathers. Birds mainly preen each other's head and neck. Known as allopreening, this behaviour is almost always done between birds of a mated pair.
Breeding
Western jackdaws become sexually mature in their second year. Genetic analysis of pairs and offspring shows no evidence of extra-pair copulation and there is little evidence for couple separation even after multiple instances of reproductive failure. Some pairs do separate in the first few months, but almost all pairings of over six months' duration are lifelong, ending only when a partner dies. Widowed or separated birds fare badly, often being ousted from nests or territories and unable to rear broods alone.
Western jackdaws usually breed in colonies with pairs collaborating to find a nest site, which they then defend from other pairs and predators during most of the year. They nest in cavities in trees or cliffs, in ruined or occupied buildings and in chimneys, the common feature being a sheltered site for the nest. The availability of suitable sites influences their presence in a locale. They may also use church steeples for nesting, a fact reported in verse by 18th century English poet William Cowper:
Nest platforms can attain a great size. A mated pair usually constructs a nest by improving a crevice by dropping sticks into it; it is then built on top of the platform formed. This behaviour has led to the blocking of chimneys and even resulted in nests crashing down into fireplaces, sometimes with birds still on them.
In his The Natural History of Selborne, Gilbert White notes that western jackdaws used to nest in crevices beneath the lintels of Stonehenge, and describes an example of the bird using a rabbit burrow for nesting. The species has been recorded outcompeting the tawny owl (Strix aluco) for nest sites in the Netherlands. They can take over old nest sites of the black woodpecker (Dryocopus martius) and stock dove (Columba oenas). Breeding colonies may also edge out those of the red-billed chough, but in turn be ousted by larger corvids such as the carrion crow, rook or magpie.
Nests are lined with hair, wool, dead grass and many other materials. The eggs are a lighter colour than those of other corvids, being smooth, a glossy pale blue or blue-green with darker speckles ranging from dark brown to olive or grey-violet. Egg size and weight varies slightly between subspecies; those of subspecies monedula average and in weight, those of subspecies soemmerringii in size and in weight, and those of subspecies spermologus in size and in weight. Clutches usually contain 4 or 5 eggs, although a Slovakian study found clutch sizes ranging from 2 to 9 eggs. The eggs are incubated by the female for 17–18 days until hatching as naked altricial chicks, which are completely dependent on the adults for food. They fledge after 28–35 days, and the parents continue to feed them for another four weeks or so.
Western jackdaws hatch asynchronously and incubation begins before clutch completion, which often leads to the death of the last-hatched young. If the supply of food is low, parental investment in the brood is kept to a minimum as little energy is wasted on feeding a chick that is unlikely to survive. Replacement clutches are very rarely laid in the event of clutch failure.
The great spotted cuckoo (Clamator glandarius) has been recorded as a brood parasite of the western jackdaw, depositing its eggs in their nests in Spain and Israel. Nest robbers include the common raven in Spain, tawny owl, and least weasel (Mustela nivalis) in England, and brown rat (Rattus norvegicus) in Finland. The European pine marten (Martes martes) raids isolated nests in Sweden but is less successful when nests are part of a colony.
Feeding
Foraging takes place mostly on the ground in open areas and to some extent in trees. Landfill sites, bins, streets, and gardens are also visited, more often early in the morning when there are fewer people about. Various feeding methods are employed, such as jumping, pecking, clod-turning and scattering, probing the soil, and occasionally, digging. Flies around cow pats are caught by jumping from the ground or at times by dropping vertically from a few metres onto the cow pat. Earthworms are not usually extracted from the ground by western jackdaws but are eaten from freshly ploughed soil. Jackdaws will ride on the backs of sheep and other mammals, seeking ticks as well as actively gathering wool or hair for nests, and will catch flying ants in flight. Compared with other corvids, the western jackdaw spends more time exploring and turning over objects with its bill; it also has a straighter and less downturned bill and increased binocular vision which are advantageous for this foraging strategy.
The western jackdaw tends to feed on small invertebrates up to in length that are found above ground, including various species of beetle (particularly cockchafers of the genus Melolontha, and weevil larvae and pupae.), Diptera, and Lepidoptera species, as well as snails and spiders. Also eaten are small rodents, bats, the eggs and chicks of birds, and carrion such as roadkill. Vegetable items consumed include farm grains (barley, wheat and oats), weed seeds, elderberries, acorns, and various cultivated fruits. Examination of the gizzards of western jackdaws shot in Cyprus in spring and summer revealed a diet of cereals (predominantly wheat) and insects (notably cicadas and beetles). The diet averages 84% plant material except when breeding, when the main food source is insects. A study in southern Spain examining western jackdaw pellets found that they contained significant amounts of silicaceous and calcareous grit to aid digestion of vegetable food and supply dietary calcium.
Opportunistic and highly adaptable, the western jackdaw varies its diet markedly depending on available food sources. They have been recorded taking eggs and nestlings from the nests of the skylark (Alauda arvensis), Manx shearwater (Puffinus puffinus), razorbill (Alca torda), common guillemot (Uria aalge), grey heron (Ardea cinerea), rock dove (Columba livia), and Eurasian collared dove (Streptopelia decaocto). A field study of a large city dump on the outskirts of León in northwestern Spain showed that western jackdaws forage there in the early morning and at dusk, and engage in some degree of kleptoparasitism. The saker falcon (Falco cherrug) has been reported stealing food from western jackdaws on powerlines in Vojvodina in Serbia.
Western jackdaws practice active food sharing – where the initiative for the transfer lies with the donor – with a number of individuals, regardless of sex or kinship. They also share more of a preferred food than a less preferred food. The active giving of food by most birds is found mainly in the context of parental care and courtship. Western jackdaws show much higher levels of active giving than has been documented for other species, including chimpanzees. The function of this behaviour is not fully understood, though it has been found to be detached from nutrition and compatible with hypotheses of mutualism, reciprocity and harassment avoidance. It has also been proposed that food sharing may be motivated by prestige enhancement.
Parasites and diseases
Western jackdaws have learned to peck open the foil caps of milk bottles left on the doorsteps after delivery by the milkman. The bacterium Campylobacter jejuni has been isolated from their beaks and cloacae so milk can become contaminated as they drink. This activity was linked to cases of Campylobacter gastroenteritis in Gateshead in northeast England and led the Department of Health to suggest that milk from bottles which had been pecked open should be discarded. It was recommended that steps be taken to prevent birds from pecking open bottles in the future.
An outbreak of a gastrointestinal illness in Spain which was causing mortalities in humans has been linked to western jackdaws. During a post-mortem on an affected bird, a polyomavirus was isolated from the spleen. The illness appeared to be a co-infection of this with Salmonella and the virus has been provisionally named the crow polyomavirus (CPyV). Segmented filamentous bacteria have been isolated from the small intestine of a western jackdaw, although their pathogenicity or role is unknown.
Pest control
The western jackdaw has been hunted as vermin, though not as heavily culled as other species of corvid. After a series of poor harvests in the early 1500s, introduced a Vermin Act in 1532 "ordeyned to dystroye Choughes (i.e. jackdaws), Crowes and Rokes" to protect grain crops from their predations. Western jackdaws were notorious as they also favoured fruit, especially cherries. This act was taken up in a piecemeal fashion, but passed the Act for the Preservation of Grayne in 1566 that was taken up with more vigour. The species was hunted for its threat to grain crops and for propensity for nesting in belfries until the mid-20th century. Particularly large numbers were culled in the county of Norfolk. Western jackdaws were also culled on game estates as they raid nests of other birds for eggs. In a 2003 dissertation on public opinion of corvids, Antonia Hereth notes that the German naturalist Alfred Brehm considered the western jackdaw to be a lovable bird, and did not describe any negative impacts of this species on agriculture.
The western jackdaw is one of a very small number of birds that it is legal to use as a decoy or to trap in a cage in the United Kingdom. The other pest species that can be controlled by trapping are the crow, jay, magpie and rook. An authorised person must comply with the requirements of the Wildlife and Countryside Act 1981 and does not need to show that the birds were a nuisance before trapping them. As of 2003 the western jackdaw was listed as a potential species for targeted hunting in the European Union Birds Directive, and hunting has been encouraged by German hunting associations. Permission to shoot western jackdaws in spring and summer exists in Cyprus as they are thought (incorrectly) to prey on gamebirds.
Cultural depictions and folklore
An ancient Greek and Roman adage runs "The swans will sing when the jackdaws are silent", meaning that educated or wise people will speak only after the foolish have become quiet. In Ancient Greek folklore, a jackdaw can be caught with a dish of oil. A narcissistic creature, it falls in while looking at its own reflection. The mythical Princess Arne Sithonis was bribed with gold by King Minos of Crete, and was punished by the gods for her greed by being transformed into an equally avaricious jackdaw, who still seeks shiny things. The Roman poet Ovid described jackdaws as harbingers of rain in his poetic work Amores. Pliny notes how the Thessalians, Illyrians, and Lemnians cherished jackdaws for destroying grasshoppers' eggs. The Veneti are fabled to have bribed the jackdaws to spare their crops.
In some cultures, a jackdaw on the roof is said to predict a new arrival; alternatively, a jackdaw settling on the roof of a house or flying down a chimney is an omen of death, and coming across one is considered a bad omen. A jackdaw standing on the vanes of a cathedral tower is said to foretell rain. The 12th-century historian William of Malmesbury records the story of a woman who, upon hearing a jackdaw chattering "more loudly than usual," grew pale and became fearful of suffering a "dreadful calamity", and that "while yet speaking, the messenger of her misfortunes arrived". Czech superstition formerly held that if jackdaws are seen quarreling, war will follow, and that jackdaws will not build nests at Sázava after being banished by Saint Procopius.
The jackdaw was considered sacred in Welsh folklore as it nested in church steeples – it was shunned by the Devil because of its choice of residence. Nineteenth century belief in the Fens held that seeing a jackdaw on the way to a wedding was a good omen for a bride.
The jackdaw is featured on the Ukrainian town of Halych's ancient coat of arms, the town's name allegedly being derived from the East Slavic word for the bird. In The Book of Laughter and Forgetting (1979), Milan Kundera notes that Franz Kafka's father Hermann had a sign in front of his shop with a jackdaw painted next to his name, since "kavka" means jackdaw in Czech.
In the video game Assassin's Creed IV: Black Flag, the main character's ship is named the Jackdaw.
| Biology and health sciences | Corvoidea | Animals |
213073 | https://en.wikipedia.org/wiki/House%20crow | House crow | The house crow (Corvus splendens), also known as the Indian, greynecked, Ceylon or Colombo crow, is a common bird of the crow family that is of Asian origin but now found in many parts of the world, where they arrived assisted by shipping. It is between the jackdaw and the carrion crow in size ( in length) but is slimmer than either. The forehead, crown, throat and upper breast are a richly glossed black, whilst the neck and breast are a lighter grey-brown in colour. The wings, tail and legs are black. There are regional variations in the thickness of the bill and the depth of colour in areas of the plumage.
Taxonomy
The nominate race C. s. splendens is found in Pakistan, India, Nepal and Bangladesh and has a grey neck collar. The subspecies C. s. zugmayeri is found in the dry parts of South Asia and Iran and has a very pale neck collar. The subspecies C. s. protegatus is found in southern India, the Maldives (sometimes separated as maledivicus) and Sri Lanka and is darker grey. C. s. insolens, found in Myanmar, is the darkest form and lacks the grey collar.
Distribution and habitat
It has a widespread distribution in southern Asia, being native to Nepal, Bangladesh, India, Pakistan, Sri Lanka, Maldives and Laccadive Islands, southern Myanmar, south Thailand, and coastal southern Iran. It was introduced to East Africa around Zanzibar (about 1897) and Port Sudan. It arrived in Australia via ship but has up to now been exterminated. Recently, it has made its arrival in Europe and has been breeding in the Dutch harbour town Hook of Holland since 1998.
A population between 200 and 400 birds has been present in Sham Shui Po, New Kowloon, Hong Kong, in particular Lai Kok Estate and Sham Shui Po Park, as well as Kowloon Tsai Park in Kowloon Tsai. An individual has been present in Cork Harbour on the south coast of Ireland since early September 2010.
In the New World, a small population of house crows is established in the area around St. Petersburg, Florida. As of April 2009, the introduced population of house crows in Socotra Island, Yemen, has been eradicated to avoid negative ecological effects to endemic species of Socotra.
It is associated with human settlements throughout its range, from small villages to large cities. In Singapore, there was a density of 190 birds/km2 in 2001 with efforts to suppress the population in planning.
Due to a human population explosion in the areas it inhabits, this species has also proportionately multiplied. Being an omnivorous scavenger has enabled it to thrive in such circumstances.
Invasive species
The invasive potential for the species is great all over the tropics. This species can make use of resources with great flexibility and appears to be associated with humans, and no populations are known to exist independently of humans. Based on niche-modelling it is estimated that it can also persist in some coastal areas in southern and northwestern Europe. In 2016, house crow was added to the List of Invasive Alien Species of Union concern. The breeding population of this species is reported from Netherlands (Hoek van Holland), while individuals were seen in Belgium, Hungary, Poland and Spain.
Behaviour
Diet
House crows feed largely on refuse around human habitations, small reptiles and mammals, and other animals such as insects and other small invertebrates, eggs, nestlings, grain and fruits. House crows have also been observed swooping down from the air and snatching baby squirrels. Most food is taken from the ground, but also from trees as opportunity arises. They are highly opportunistic birds and given their omnivorous diet, they can survive on nearly anything that is edible. These birds can be seen near marketplaces and garbage dumps, foraging for scraps. They have also been observed to eat sand after feeding on carcasses.
Nesting
At least some trees in the local environment seem to be necessary for successful breeding although house crows occasionally nest on telephone towers. It lays 3–5 eggs in a typical stick nest, and occasionally there are several nests in the same tree. In South Asia they are parasitized by the Asian koel. Peak breeding in India as well as Peninsular Malaysia is from April to July. Large trees with big crowns are preferred for nesting.
Roosting
House crows roost communally near human habitations and often over busy streets. A study in Singapore found that the preferred roost sites were in well-lit areas with a lot of human activity, close to food sources and in tall trees with dense crowns that were separated from other trees. The roost sites were often enclosed by tall buildings. Before flying into roost trees, crows make pre-roosting aggregations perched on TV antennas, roof tops, wayside trees, open fields, and feed or preen during this time.
Voice
The voice is a harsh kaaw-kaaw.
Relationship to humans
It is suspected that paramyxoviruses, such as PMV 1 that causes of Newcastle disease may be spread by Corvus splendens. Outbreaks of Newcastle disease in India were often preceded by mortality in crows. They have also been found to carry Cryptococcus neoformans, which can cause cryptococcosis in humans.
House crows in Tanzania curiously showed an absence of blood parasites, although some species such as Trypanosoma corvi have been first described from this species. Pathologist T.R. Lewis expressed surprise at the numbers of haematozoa present in the blood of house crows from Calcutta.
Gallery
| Biology and health sciences | Corvoidea | Animals |
213108 | https://en.wikipedia.org/wiki/American%20crow | American crow | The American crow (Corvus brachyrhynchos) is a large passerine bird species of the family Corvidae. It is a common bird found throughout much of North America. American crows are the New World counterpart to the carrion crow and the hooded crow of Eurasia; they all occupy the same ecological niche. Although the American crow and the hooded crow are very similar in size, structure and behavior, their calls and visual appearance are different.
From beak to tail, an American crow measures , almost half of which is tail. Its wingspan is . Mass varies from about , with males tending to be larger than females. Plumage is all black, with iridescent feathers. It looks much like other all-black corvids. They are very intelligent, and adaptable to human environments. The most usual call is They can be distinguished from the common raven (C. corax) because American crows are smaller and the beak is slightly less pronounced; from the fish crow (C. ossifragus) because American crows do not hunch and fluff their throat feathers when they call; and from the carrion crow (C. corone) by size, as the carrion crow is larger and of a stockier build.
American crows are common, widespread, and susceptible to the West Nile virus, making them useful as a bioindicator to track the virus's spread. Direct transmission of the virus from crows to humans is impossible.
Taxonomy and systematics
The American crow was described by German ornithologist Christian Ludwig Brehm in 1822. Its scientific name means literally 'short-billed crow', from Ancient Greek ('short-') and ('billed').
A 2012 genetic analysis of the genus Corvus by Knud A Jønsson and colleagues using nuclear and mitochondrial DNA calculated that the American crow diverged from a lineage that gave rise to the collared, carrion and hooded crows around 5 million years ago.
"American crow" has been designated the official name by the International Ornithologists' Union (IOC).
Subspecies
The number of subspecies varies by authority, ranging between three and five. The unclear taxonomy of the northwestern crow, previously its own species, has complicated subspecies determinations. Subspecies differ in bill proportion and form a rough NE–SW clinal in size across North America. Birds are smallest in the far west and on the southern coast.
C. b. brachyrhynchos (Brehm, 1822) – eastern crow: northeastern United States, eastern Canada and surroundings. The nominate subspecies, and largest.
C. b. hesperis (Ridgway, 1887) – western crow: western North America except the Arctic north, the Pacific Northwest and the extreme south. Smaller overall with a proportionally more slender bill and low-pitched voice.
C. b. caurinus (Baird, 1858) – northwestern crow: of the Pacific temperate rain forests was formerly considered a distinct species as C. caurinus, averaging smaller in size than other American crows with a distinctly hoarser call. It forms a hybrid swarm with American crow (sensu stricto) in coastal Washington and British Columbia. The American Ornithological Society lumped the northwestern crow with the American crow in 2020. It is now considered a geographic variation within C. b. hesperis.
C. b. pascuus (Coues, 1899) – Florida crow: Florida. Mid-sized, short-winged, but decidedly long bill and legs.
C. b. paulus (Howell, 1913) – southern crow: southern United States. Smaller overall, bill also small.
Description
The American crow is a large, distinctive bird with iridescent black feathers all over. Its legs, feet and bill are also black. They measure in length, of which the tail makes up about 40%. The wing chord is , with the wingspan ranging from . The bill length can be from , varying strongly according to location. The tarsus is and the tail is . The body mass can vary from . Males tend to be larger than females.
The most usual call is a loud, short, and rapid . Usually the birds thrust their heads up and down as they utter this call. American crows can also produce a wide variety of sounds and sometimes mimic noises made by other animals, including other birds, such as barred owls.
Visual differentiation from the fish crow (C. ossifragus) is extremely difficult and often inaccurate. Nonetheless, differences apart from size do exist. Fish crows tend to have more slender bills and feet. There may also be a small sharp hook at the end of the fish crow's upper bill. Fish crows also appear as if they have shorter legs when walking. More dramatically, when calling, fish crows tend to hunch and fluff their throat feathers.
If seen flying at a distance from where size estimates are unreliable, the distinctly larger common ravens (C. corax) can be distinguished by their almost lozenge-shaped tail and their larger-looking heads. They also fluff their throat feathers when calling – like fish crows, only more so. Ravens also soar for extended periods, unlike crows, which rarely fly more than a few seconds without flapping their wings.
Crows have been noted to be intelligent. They have the same brain-weight-to-body ratio as humans. This has led to some studies that have identified that crows are self-aware and that young crows take time to learn from tolerant parents. While a human has a neocortex, the crow has a different area in their brain that is equally complex.
The average lifespan of the American crow in the wild is 7–8 years. Captive birds are known to have lived up to 30 years.
Distribution and habitat
The range of the American crow now extends from the Pacific Ocean to the Atlantic Ocean in Canada, on the French islands of Saint-Pierre and Miquelon, south through the United States, and into northern Mexico. The increase in trees throughout the Great Plains during the past century due to fire suppression and tree planting facilitated range expansions of the American crow as well as range expansions of many other species of birds. Virtually all types of country from wilderness, farmland, parks, open woodland to towns and major cities are inhabited; it is absent only from tundra habitat, where it is replaced by the common raven. This crow is a permanent resident in most of the US, but most Canadian birds migrate some distances southward in winter. Outside of the nesting season these birds often gather in large (thousands or even millions) communal roosts at night.
The American crow was recorded in Bermuda from 1876 onwards.
Behavior and ecology
Studying the behavior of American crows is laborious due to the difficulty in catching them to band them, let alone catching them again. Thus much of their behavior, including daily routine, migration, molting, survivorship, age of first breeding, nestling development, and the nature of nesting helpers, remains poorly studied.
Diet
The American crow is omnivorous. It will feed on invertebrates of all types, carrion, scraps of human food, fruits, nuts such as walnuts and almonds, seeds, eggs and nestlings, stranded fish on the shore and various grains. American crows are active hunters and will prey on mice, young rabbits, frogs, and other small animals. In the winter and autumn, the diet of American crows is more dependent on nuts and acorns. Occasionally, they will visit bird feeders. The American crow is one of only a few species of bird that has been observed modifying and using tools to obtain food.
Like most crows, they will scavenge at landfills, scattering garbage in the process. Where available, corn, wheat and other crops are a favorite food. These habits have historically caused the American crow to be considered a nuisance. However, it is suspected that the harm to crops is offset by the service the American crow provides by eating insect pests.
Reproduction
American crows are socially monogamous cooperative breeding birds. Mated pairs form large families of up to 15 individuals from several breeding seasons that remain together for many years. Offspring from a previous nesting season will usually remain with the family to assist in rearing new nestlings. American crows do not reach breeding age for at least two years. Most do not leave the nest to breed for four to five years.
The nesting season starts early, with some birds incubating eggs by early April. American crows build bulky stick nests, nearly always in trees but sometimes also in large bushes and, very rarely, on the ground. They will nest in a wide variety of trees, including large conifers, although oaks are most often used. Three to six eggs are laid and incubated for 18 days. The young are usually fledged by about 36 days after hatching. Predation primarily occurs at the nest site and eggs and nestlings are frequently eaten by snakes, raccoons, ravens and domestic cats. Adults are less frequently predated, but face potential attack from great horned owls, red-tailed hawks, peregrine falcons and eagles. They may be attacked by predators such as coyotes or bobcats at carrion when incautious, although this is even rarer.
Although matings between relatives can cause inbreeding depression, that is, negative effects on offspring fitness, a number of species never-the-less prefer inbreeding. The American crow experiences inbreeding depression (as measured by fledgling survival) but still prefers inbreeding.
West Nile virus
American crows succumb easily to West Nile virus infection. This was originally a mosquito-borne African virus causing encephalitis in humans and livestock since about 1000 AD, and was accidentally introduced to North America in 1999, apparently by an infected air traveller who got bitten by a mosquito after arrival. It is estimated that the American crow population has dropped by up to 45% since 1999. Despite this decline, the crow is considered a species of least concern.
The disease runs most rampant in the subtropical conditions which encourage reproduction of its mosquito vectors among which Culex tarsalis is most significant. Mortality rates appear to be higher than those in other birds, causing local population losses of up to 72% in a single season. Because of this, American crows are a sentinel species indicating the presence of West Nile virus in an area. Crows cannot transmit the virus to humans directly.
Intelligence
American crows, like other corvids, are highly cunning and inquisitive. They are able to steal food from other species, often in creative ways. One example shows a group of crows stealing a fish from a Northern river otter: one bird pecked the otter's tail to distract it while other birds swooped in and stole the fish. They are able to use and modify tools.
Relationship with humans
Crows are used as a motif in some human cultures, often associated with death, thieves, graveyards, bad luck, and other negative connotations.
Conversely, they are seen by some neo-pagan and indigenous cultures as signs of good luck, or even signs of certain gods, such as Apollo, Odin, and others.
Status and management
The intelligence and adaptability of the American crow has insulated it from threats, and it is instead considered an agricultural pest. In 2012, BirdLife International estimated the American crow population to be around 31 million. The large population and vast range result in the least concern status for the American crow, meaning that the species is not threatened with extinction.
Crows have been killed in large numbers by humans, both for recreation and as part of organized campaigns of extermination. In Canada, American crows have no protections, aside from Quebec, which bans their hunting during the nesting season. Laws on their hunting vary throughout the United States. New Jersey allows for a limited hunting season, unless they are agricultural pests, in which case they may be killed. Oklahoma allows hunting even during the nesting season. In the first half of the 20th century, state-sponsored campaigns dynamited roosting areas, taking large numbers of crows. A campaign in Oklahoma from 1934 to 1945 dynamited 3.8 million birds. The effect on populations was negligible and damage to agricultural crops did not decrease, and thus the campaign was halted due to its ineffectiveness. In a study taking data from 1917 to 1999, intentional killings were the overwhelming cause of death for crows, accounting for 68% of all recovered bird bands.
Non-deadly methods of managing crows are varied but usually limited in their effectiveness. High value crops may be netted, but this is cost prohibitive for most other crops. Frightening may be used to disperse crows, including loud noises from guns, fake hawks flown from balloons, fake owls that move with the wind, strips of reflective tape on fences, or recordings of crow distress calls. Poisoned baits are of limited effectiveness, as only the most toxic baits work, and those are generally unacceptable for use. Crows quickly learn to avoid the less-toxic baits, as the baits make crows sick. The actual effect of crows on agriculture has been poorly studied. There is some suggestion that they may be a benefit to farmers, by eating insect pests and chasing off livestock predators like hawks.
| Biology and health sciences | Corvoidea | Animals |
213403 | https://en.wikipedia.org/wiki/Primula | Primula | Primula () is a genus of herbaceous flowering plants in the family Primulaceae. They include the primrose (P. vulgaris), a familiar wildflower of banks and verges. Other common species are P. auricula (auricula), P. veris (cowslip), and P. elatior (oxlip). These species and many others are valued for their ornamental flowers. They have been extensively cultivated and hybridised (in the case of the primrose, for many hundreds of years). Primula are native to the temperate Northern Hemisphere, south into tropical mountains in Ethiopia, Indonesia, and New Guinea, and in temperate southern South America. Almost half of the known species are from the Himalayas.
Primula has over 500 species in traditional treatments, and more if certain related genera are included within its circumscription.
Description
Primula is a complex and varied genus, with a range of habitats from alpine slopes to boggy meadows. Plants bloom mostly during the spring, with flowers often appearing in spherical umbels on stout stems arising from basal rosettes of leaves; their flowers can be purple, yellow, red, pink, blue, or white. Some species show a white mealy bloom (farina) on various parts of the plant. Many species are adapted to alpine climates.
Taxonomy
Primula was known at least as early as the mediaeval herbalists, although first formally described as a genus by Linnaeus in 1753, and later in 1754 in his Flora Anglica. Linnaeus described seven species of Primula. One of its earliest scientific treatments was that of Charles Darwin study of heterostyly in 1877 (The different forms of flowers on plants of the same species). Since then, heterostyly (and homostyly) have remained important considerations in the taxonomic classification of Primula. Primula is a member of the Primulaceae family. The most complete treatment of the family, with nearly 1000 species arranged into 22 genera, was by Pax and Knuth in 1905.
Phylogeny
Primula is the largest genus in the family Primulaceae, within which it is placed in the subfamily Primuloideae, being the nominative genus.
The position of Primula within the family and its relationship to other genera is shown in this cladogram:
Classification
The genus Dodecatheon originated from within Primula, its species are now included in Primula.
Sections of genus Primula
The classification of the genus Primula has been investigated by botanists for over a century. As the genus is both large and diverse (with about 430–500 species), botanists have organized the species in various sub-generic groups. The most common is division into a series of thirty sections. Some of these sections (e.g. Vernales, Auricula) contain many species; others contain only one.
Amethystina
Auricula
Bullatae
Candelabra
Capitatae
Carolinella
Cortusoides
Cuneifolia
Denticulata
Dryadifolia
Farinosae
Floribundae
Grandis
Malacoidea
Malvacea
Minutissimae
Muscaroides
Nivales
Obconica
Parryi
Petiolares
Pinnatae
Pycnoloba
Reinii
Rotundifolia
Sikkimensis
Sinenses
Soldanelloideae
Souliei
Vernales
Selected species
Species include:
Primula alcalina (bluedome primrose)
Primula algida
Primula aliciae
Primula allionii (Allioni's primrose)
Primula alpicola (moonlight primrose)
Primula amethystina
Primula angustifolia (alpine primrose)
Primula anisodora (anise primrose)
Primula anvilensis (boreal primrose)
Primula appenina
Primula arunachalensis
Primula atrodentata
Primula aurantiaca (primevère à fleurs oranges)
Primula aureata
Primula auricula (auricula, bear's ear)
Primula auriculata
Primula bathangensis
Primula beesiana (candelabra primrose)
Primula bella
Primula bellidifolia
Primula bergenioides
Primula bhutanica
Primula blattariformis
Primula boothii
Primula borealis (northern primrose)
Primula bracteosa
Primula bulleyana (candelabra primrose)
Primula calderiana
Primula calliantha
Primula calyptrata
Primula capillaris (Ruby Mountain primrose)
Primula capitata
Primula capitata ssp. mooreana
Primula capitata ssp. crispata
Primula capitellata
Primula carniolica (Carniolan primrose)
Primula caveana
Primula cawdoriana
Primula cernua
Primula chionantha
Primula chumbiensis
Primula chungensis
Primula chasmophila
Primula clarkei
Primula clusiana
Primula clutterbuckii
Primula cockburniana
Primula concholoba
Primula concinna
Primula conspersa
Primula cortusoides
Primula cuneifolia (wedgeleaf primrose, pixie-eye primrose)
Primula cusickiana (Cusick's primrose)
Primula daonensis
Primula darialica
Primula davidii
Primula deflexa
Primula denticulata (drumstick primrose, Himalayan primrose)
Primula deorum (Rila primrose, Rila cowslip, God's cowslip)
Primula deuteronana
Primula dryadifolia
Primula edgeworthii
Primula egaliksensis (Greenland primrose)
Primula elatior (oxlip, true oxlip, oxslip)
Primula elizabethiae
Primula elongata
Primula erythrocarpa
Primula farinosa (birdseye primrose)
Primula farreriana
Primula fedschenkoi
Primula fenghwaiana
Primula filchnerae
Primula filipes
Primula firmipes
Primula fistulosa
Primula flaccida
Primula floribunda
Primula florindae (Himalayan cowslip, Tibetan cowslip)
Primula forrestii
Primula frondosa
Primula gambeliana
Primula gemmifera
Primula geraniifolia
Primula giraldiana
Primula glabra
Primula glaucescens
Primula glomerata
Primula glutinosa
Primula gracilenta
Primula gracillipes
Primula griffithii
Primula halleri (long-flowered primrose, Haller's primrose)
Primula helodoxa
Primula hendersonii (broad-leaved shootingstar, Henderson's shootingstar, mosquito-bill, sailor caps)
Primula heucherifolia
Primula hirsuta (stinking primrose)
Primula hookeri
Primula incana (silvery primrose, mealy primrose)
Primula interjacens
Primula involucrata
Primula ioessa
Primula irregularis
Primula japonica (Japanese primrose, Japanese cowslip)
Primula jesoana
Primula jigmediana
Primula juliae (Julia's primrose, purple primrose)
Primula kingii
Primula kisoana
Primula kitaibeliana (Kitaibel's primrose)
Primula klattii
Primula laurentiana (birdseye primrose)
Primula lihengiana
Primula listeri
Primula longiscapa
Primula lutea
Primula luteola
Primula macrophylla (largeleaf primrose)
Primula magellanica
Primula malacoides (fairy primrose, baby primrose)
Primula malvacea
Primula marginata (marginate primrose)
Primula megaseifolia
Primula melanantha
Primula melanops
Primula minima (dwarf primrose)
Primula minor
Primula mistassinica (Mistassini primrose)
Primula modesta
Primula mollis
Primula moupinensis
Primula munroi
Primula muscarioides
Primula nipponica
Primula nivalis (snowy primrose)
Primula obconica (poison primrose, German primrose)
Primula orbicularis
Primula palinuri
Primula parryi (Parry's primrose)
Primula pedemontana
Primula petelotii
Primula petiolaris
Primula pinnatifida
Primula poissonii
Primula polyneura
Primula prenantha
Primula prolifera (glory-of-the-marsh)
Primula pulchella
Primula pulverulenta (mealy cowslip)
Primula pumila (arctic primrose)
Primula reidii
Primula reinii
Primula renifolia
Primula reptans
Primula reticulata
Primula rockii
Primula rosea (Himalayan meadow primrose)
Primula rotundifolia
Primula rusbyi (Rusby's primrose)
Primula sapphirina
Primula saxatilis (rock primrose)
Primula scandinavica (Scandinavian primrose)
Primula scapigera
Primula scotica (Scottish primrose)
Primula secundiflora
Primula septemloba
Primula serratifolia
Primula sheriffii
Primula sieboldii (Japanese primrose)
Primula sikkimensis (Sikkim cowslip)
Primula sinensis (syn. P. praenitens)
Primula sinolisteri
Primula sinomollis
Primula soldanelloides
Primula sonchifolia
Primula souliei
Primula spectabilis
Primula specuicola (alcove primrose, cave-dwelling primrose)
Primula stricta (coastal primrose, strict primrose)
Primula strumosa
Primula suffrutescens (Sierra primrose)
Primula szechuanica
Primula takedana
Primula tanneri
Primula tardiflora
Primula tenella
Primula tenuiloba
Primula tenuituba
Primula tibetica
Primula tschuktschorum (Chukchi primrose)
Primula tyrolensis
Primula vaginata
Primula valentiniana
Primula veris (cowslip)
Primula verticillata (yellow primrose)
Primula vialii (wayside primrose, pagoda primrose, orchid primrose, poker primrose)
Primula villosa
Primula vulgaris (primrose)
Primula waltonii
Primula watsonii
Primula warshenewskiana
Primula whitei
Primula wilsonii
Primula wollastonii (Wollaston's primrose)
Primula wulfeniana (Wulfen's primrose)
Primula xanthopa
Primula yunnanensis
Primula yuparensis
Etymology
The word primula is the Latin feminine diminutive of primus, meaning first (prime), applied to flowers that are among the first to open in spring.
Distribution and habitat
Although there are over 400 species of Primula, about 75% are found in the eastern Himalayan mountain chain and western China (Yunnan Province), constituting a centre of diversity. Other centres of diversity are a western Asian centre (Caucasus, European mountain ranges from the Pyrenees, through the Alps to the Carpathian Mountains), mountains of East Asia and those of western North America. Primula is found in mountainous or higher latitude zones of North America, Europe, and Asia, with extension into South America, Africa (mountains of Ethiopia) and tropical Asia (islands of Java and Sumatra). About 25 species occur in North America (represented in five sections).
Primula is found in the humid and moderate climate regions of the Northern Hemisphere, predominantly in the forest belt, plain meadows, Alpine lawns, and nival and meadow tundras.
Ecology
Primulas are used as a food plant by the Duke of Burgundy butterfly.
Cultivation
Primula species have been extensively cultivated and hybridised, mainly derived from P. elatior, P. juliae, P. veris and P. vulgaris. Polyanthus (Primula × polyantha) is one such group of plants, which has produced a large variety of strains in all colours, usually grown as annuals or biennials and available as seeds or young plants.
Another huge range of cultivars, known as auriculas, are derived from crosses between P. auricula and P. hirsuta (among others). Specialist nurseries and auricula societies support the growing and showing of these choice strains.
AGM cultivars
The following hybrid varieties and cultivars have gained the Royal Horticultural Society's Award of Garden Merit:-
'Broadwell Milkmaid' (auri)
Charisma series (prim)
'Clarence Elliott' (auri)
Crescendo Series (poly)
'Crescendo Blue Shades'
'Crescendo Bright Red'
'Crescendo Pink and Rose Shades'
Danova Series (prim)
'Danova Rose'
'Francisca' (poly)
'Guinevere' (poly)
Primula × loiseleurii 'Aire Mist' (auri)
Primula × pubescens (auri)
'Tony' (auri)
'Wanda' (prim)
Hybrids
Primula × kewensis=P. floribunda × P. verticillata (Kew primrose)
Primula × polyantha=P. veris × P. vulgaris (false oxlip, polyanthus primula)
Primula × pubescens=P. hirsuta × P. auricula
| Biology and health sciences | Ericales | Plants |
213407 | https://en.wikipedia.org/wiki/Jay | Jay | A jay is a member of a number of species of medium-sized, usually colorful and noisy, passerine birds in the crow family, Corvidae. The evolutionary relationships between the jays and the magpies are rather complex. For example, the Eurasian magpie seems more closely related to the Eurasian jay than to the East Asian blue and green magpies, whereas the blue jay is not closely related to either. The Eurasian jay distributes oak acorns, contributing to the growth of oak woodlands over time.
Systematics and species
Jays are not a monophyletic group. Anatomical and molecular evidence indicates they can be divided into a New World and an Old World lineage (the latter including the ground jays and the piapiac), while the grey jays of the genus Perisoreus form a group of their own. The black magpies, formerly believed to be related to jays, are classified as treepies.
Old World ("brown") jays
Grey jays
New World jays
In culture
Slang
The word jay has an archaic meaning in American slang meaning a person who chatters impertinently.
The term jaywalking was coined in the first decade of the 1900s to label persons crossing a busy street carelessly and becoming a traffic hazard. The term began to imply recklessness or impertinent behavior as the convention became established.
In January 2014, Canadian author Robert Joseph Greene embarked on a lobbying campaign among ornithologists in Europe and North America to get Merriam-Websters Dictionary to have a "Jabber of Jays" as an official term under bird groups.
| Biology and health sciences | Corvoidea | Animals |
213429 | https://en.wikipedia.org/wiki/Accentor | Accentor | The accentors are a genus of birds in the family Prunellidae, which is endemic to the Old World. This small group of closely related passerines are all in the genus Prunella. All but the dunnock and the Japanese accentor are inhabitants of the mountainous regions of Europe and Asia; these two also occur in lowland areas, as does the Siberian accentor in the far north of Siberia. These birds are not strongly migratory, but they will leave the coldest parts of their range in winter and make altitudinal movements.
Taxonomy and etymology
The genus Prunella was introduced by the French ornithologist Louis Vieillot in 1816 with the dunnock (Prunella modularis) as the type species. Although the genus is usually used for all the accentors, the alpine accentor and Altai accentor are sometimes separated into the genus Laiscopus.
Harrison used the group name dunnock for all of the species, not just Prunella modularis (thus e.g. Japanese dunnock for P. rubida); this usage is based on the oldest known name for any of the species (old English dun-, brown, + -ock, small: "little brown bird"). Accentor derives from the old scientific name for the Alpine accentor (Accentor collaris). It comes from Late Latin, meaning "sing with another" (ad + cantor). The genus name Prunella is from the German Braunelle, "dunnock", a diminutive of braun, "brown".
Description
These are small, fairly drab species superficially similar, but not closely related to, sparrows. They are 14 to 18 cm in length, and weigh between 25 and 35 g. However, accentors have thin sharp bills, reflecting their diet of ground-dwelling insects in summer, augmented with small seeds and berries in winter. They may also swallow grit and sand to help their stomach break up these seeds.
Most of the species live together in flocks. The dunnock is an exception since it prefers to be solitary except when feeding. The dunnock also earned a nickname of "shuffle-wing" since it most strongly displays the characteristic wing flicks used during courtship and other displays.
Accentors may have one to three broods a year. Courtship consists of a great deal of song from the males, which may include short lark-like song flights to attract a mate. In most species, the male and female share in the nest making, with the dunnocks again being an exception; their males have no part in nest building or incubation. They build neat cup nests and lay 3–6 unspotted green or blue eggs. The eggs are incubated for around 12 days. The young are fed by both parents and take an additional 12 days or so to fledge.
Habitat
Their typical habitat is mountainous regions in an area far above the tree-line, but below the snow-line. The Alpine accentor has been observed at nearly altitude in the Himalaya and the Altai and Robin accentors at ; however, most accentors breed in scrub vegetation at lower levels. Most species are altitudinal migrants, descending to lower altitudes to spend the winter, but some are long-distance migrants, most notably the Siberian accentor, which migrates 3,000–5,000 km from northern Siberia to China and Korea. Accentors spend the majority of their time in the undergrowth and even when flushed, stay low to the ground until reaching cover.
Species list
Twelve species are currently accepted:
Alpine accentor Prunella collaris
Altai accentor Prunella himalayana
Robin accentor Prunella rubeculoides
Rufous-breasted accentor Prunella strophiata
Siberian accentor Prunella montanella
Brown accentor Prunella fulvescens
Radde's accentor Prunella ocularis (includes the formerly accepted Arabian accentor Prunella fagani, lumped in 2018)
Black-throated accentor Prunella atrogularis
Kozlov's accentor Prunella koslowi
Dunnock Prunella modularis
Japanese accentor Prunella rubida
Maroon-backed accentor Prunella immaculata
| Biology and health sciences | Passerida | Animals |
213454 | https://en.wikipedia.org/wiki/Corvus%20%28constellation%29 | Corvus (constellation) | Corvus is a small constellation in the Southern Celestial Hemisphere. Its name means "crow" in Latin. One of the 48 constellations listed by the 2nd-century astronomer Ptolemy, it depicts a raven, a bird associated with stories about the god Apollo, perched on the back of Hydra the water snake. The four brightest stars, Gamma, Delta, Epsilon, and Beta Corvi, form a distinctive quadrilateral or cross-shape in the night sky.
With an apparent magnitude of 2.59, Gamma Corvi—also known as Gienah—is the brightest star in the constellation. It is an aging blue giant around four times as massive as the Sun. The young star Eta Corvi has been found to have two debris disks. Three star systems have exoplanets, and a fourth planetary system is unconfirmed. TV Corvi is a dwarf nova—a white dwarf and brown dwarf in very close orbit.
History and mythology
In the Babylonian star catalogues dating from at least 1100 BCE, what later became known as Corvus was called the Raven (MUL.UGA.MUSHEN). As with more familiar Classical astronomy, it was placed sitting on the tail of the Serpent (Greek Hydra). The Babylonian constellation was sacred to Adad, the god of rain and storm; in the second Millennium BCE it would have risen just before the autumnal rainy season. John H. Rogers observed that Hydra signified Ningishzida, the god of the underworld in the Babylonian compendium MUL.APIN. He proposed that Corvus and Crater (along with Hydra) were death symbols and marked the gate to the underworld. These two constellations, along with the eagle Aquila and the fish Piscis Austrinus, were introduced to the Greeks around 500 BCE; they marked the winter and summer solstices respectively. Furthermore, Hydra had been a landmark as it had straddled the celestial equator in antiquity. Corvus and Crater also featured in the iconography of Mithraism, which is thought to have been of middle-eastern origin before spreading into Ancient Greece and Rome.
Corvus is associated with the myth of Apollo and his lover Coronis the Lapith. Coronis had been unfaithful to Apollo; when he learned this information from a pure white crow (or raven in some versions), he turned its feathers black in a fit of rage. Another legend associated with Corvus is that a crow stopped on his way to fetch water for Apollo, to eat figs. Instead of telling the truth to Apollo, he lied and said that a snake, Hydra, kept him from the water, while holding a snake in his talons as proof. Apollo, realizing this was a lie, flung the crow (Corvus), cup (Crater), and snake (Hydra) into the sky. He further punished the wayward bird by ensuring it would forever be thirsty, both in real life and in the heavens, where the Cup is just out of reach.
In other cultures
In Chinese astronomy, the stars of Corvus are located within the Vermilion Bird of the South (南方朱雀, Nán Fāng Zhū Què). The four main stars depict a chariot, Zhen, which is the 28th and final lunar mansion; Alpha and Eta mark the linchpins for the wheels, and Zeta is Changsha, a coffin. In Indian astronomy, the five main stars of Corvus represent a hand or fist corresponding to the Hasta, the 13th nakshatra or lunar mansion.
Corvus was recognized as a constellation by several Polynesian cultures and used as a guide for ocean navigation.
In the Marquesas Islands, it was called Mee; in Pukapuka, it was called Te Manu, and in the Society Islands, it was called Metua-ai-papa. To Torres Strait Islanders, Corvus was the right hand (holding kupa fruit) of the huge constellation Tagai, a man fishing.
The Bororo people of Mato Grosso in central Brazil regarded the constellation as a land tortoise Geriguigui, while the Tucano people of the northwestern Amazon region saw it as an egret. To the Tupi people of São Luís Island in Brazil, Corvus might have been seen as a grill or barbecue—seychouioura, on which fish were grilled. The depiction could have also referred to the Great Square of Pegasus.
Characteristics
Covering 184 square degrees and hence 0.446% of the sky, Corvus ranks 70th of the 88 constellations in area. It is bordered by Virgo to the north and east, Hydra to the south, and Crater to the west. The three-letter abbreviation for the constellation, as adopted by the International Astronomical Union in 1922, is "Crv". The official constellation boundaries, as set by Belgian astronomer Eugène Delporte in 1930, are defined by a polygon of six segments (illustrated in infobox). In the equatorial coordinate system, the right ascension coordinates of these borders lie between and , while the declination coordinates are between −11.68° and −25.20°. Its position in the Southern Celestial Hemisphere means that the whole constellation is visible to observers south of 65°N.
Features
Stars
The German cartographer Johann Bayer used the Greek letters Alpha through Eta to label the most prominent stars in the constellation. John Flamsteed gave nine stars Flamsteed designations, while one star he designated in the neighbouring constellation Crater—31 Crateris—lay within Corvus once the constellation boundaries were established in 1930. Within the constellation's borders, there are 29 stars brighter than or equal to apparent magnitude 6.5.
Four principal stars, Delta, Gamma, Epsilon, and Beta Corvi, form a quadrilateral asterism known as "the "Spica's Spanker" or "the Sail". Although none of the stars are particularly bright, they lie in a dim area of the sky, rendering the asterism easy to distinguish in the night sky. Gamma and Delta serve as pointers toward Spica.
Also called Gienah, Gamma is the brightest star in Corvus at magnitude 2.59. Its traditional name means "wing", the star marking the left wing in Bayer's Uranometria. 154±1 light-years from Earth, it is a blue-white hued giant star of spectral type B8III that is times as massive, and 355 times as luminous as the Sun. Around million years old, it has largely exhausted its core hydrogen and begun expanding and cooling as it moves away from the main sequence. A binary star, it has a companion orange or red dwarf star of spectral type K5V to M5V that is about 0.8 times as massive as the Sun. Around 50 astronomical units distant from Gamma Corvi A, it is estimated to complete an orbit in 158 years.
Delta Corvi, traditionally called Algorab, is a double star divisible in small amateur telescopes. The primary is a blue-white star of magnitude 2.9, around 87 light-years from Earth. An enigmatic star around 2.7 times as massive as the Sun, it is more luminous (65–70 times that of the Sun) than its should be for its surface temperature of 10,400 K, and hence is either a 3.2 million year-old very young pre-main sequence star that has not settled down to a stable main sequence life stage, or a 260-million-year-old star that has begun to exhaust its core hydrogen and expand, cool and shine more brightly as it moves away from the main sequence. Its spectral type is given as A0IV, corresponding with the latter scenario. Warm circumstellar dust—by definition part of its inner stellar system—has been detected around Delta Corvi A. Delta Corvi B is an orange dwarf star of magnitude 8.51 and spectral class K, also surrounded by circumstellar dust. A post T-tauri star, it is at least 650 AU distant from its brighter companion and takes at least 9400 years to complete an orbit. Delta Corvi's common name means "the raven". It is one of two stars marking the right wing. Located 4.5 degrees northeast of Delta Corvi is Struve 1669, a binary star that is divisible into two stars 5.4" apart by small amateur telescopes, 280 light-years from Earth. The pair, both white stars, are visible to the naked eye at magnitude 5.2; the primary is of magnitude 5.9 and the secondary is of magnitude 6.0.
The raven's breast is marked by Beta Corvi (the proper name is Kraz), a star of magnitude 2.7 located 146 ± 1 light-years from Earth. Roughly 206 million years old and 3.7 ± 1 times as massive as the Sun, it has exhausted its core hydrogen and expanded and cooled to a surface temperature of around 5,100 K and is now a yellow bright giant star of spectral type G5II. It likely spent most of its existence as a blue-white main sequence star of spectral type B7V. Bearing the proper name of Minkar and marking the raven's nostril is Epsilon Corvi, located some 318 ± 5 light-years from Earth. It is a red giant of spectral type K2III that is around 54 times the Sun's radius and 930 times its luminosity. Around 4 times as massive as the Sun, it spent much of its life as a main-sequence star of spectral type B5V. Lying to the south of the quadrilateral between Beta and Epsilon Corvi is the orange-hued 6 Corvi, an ageing giant star of spectral type K1III that is around 70 times as luminous as the Sun. It is 331 ± 10 light-years away from Earth.
Named Alchiba, Alpha Corvi is a white-hued star of spectral type F1V and magnitude 4.0, 48.7 ± 0.1 light-years from Earth. It exhibits periodic changes in its spectrum over a three-day period, which suggests it is either a spectroscopic binary or (more likely) a pulsating Gamma Doradus-type variable. If the latter is the case, it is estimated to be 1.39 times as massive as the Sun. According to Bayer's atlas, it lies above the bird's beak.
Marking the raven's right wing is Eta Corvi, a yellow-white main-sequence star of type F2V that is 1.52 times as massive and 4.87 times as luminous as the Sun. It is 59 light-years distant from our Solar System. Two debris disks have been detected orbiting this star, one warm within 3.5 astronomical units and another out at ~150 astronomical units distant. Zeta Corvi marks the raven's neck. It is of apparent magnitude 5.21, separated by 7 arcseconds from the star HR 4691. Located 420 ± 10 light-years distant, it is a blue-white Be star of spectral type B8V, the presence of hydrogen emission lines in its spectrum indicating it has a circumstellar disc. These stars may be an optical double or a true multiple star system, with a separation of at least 50,000 astronomical units and the stars taking 3.5 million years to orbit each other. HR 4691 is itself double, composed of an ageing yellow-orange giant whose spectral type has been calculated at K0 or G3, and an F-type main-sequence star.
31 Crateris (which was originally placed in Crater by Flamsteed) is a 5.26 magnitude star which was once mistaken for a moon of Mercury. On 27 March 1974, the Mariner 10 mission detected emissions in the far ultraviolet from the planet (suggesting a satellite), but they were found to emanate from the star. It is in reality a remote binary star system with a hot blue-white star of spectral type B1.5V and a companion about which little is known. The two stars orbit each other every 2.9631 days. The primary is possibly a blue straggler of the Hyades group. The primary is around 15.5 times as massive as the Sun and 52262 times as luminous.
VV Corvi is a close spectroscopic binary, its two component stars orbiting each other with a period of 1.46 days. Both are yellow-white main-sequence stars of spectral type F5V, though the primary has begun expanding and cooling as it nears the end of its time on the main sequence. The mass ratio of the two stars is 0.775 ± 0.024. A tertiary companion was discovered during the Two Micron All-Sky Survey. W Corvi is an eclipsing binary that varies in brightness from apparent magnitude 11.16 to 12.5 over 9 hours. Its period has increased by 1/4 second over a century. It is an unusual system in that its two stars are very close to each other yet have different surface temperatures, and hence thermal transfer is not taking place as expected. SX Corvi is an eclipsing binary that is also a contact binary known as a W Ursae Majoris variable. The two component stars orbit closely enough to each other for mass to have been transferred between them—in this case the secondary having transferred a large amount of mass to the primary. RV Corvi is another eclipsing binary. Its brightness varies from apparent magnitude 8.6 to 9.16 over 18 hours. The system is composed of stars of spectral types F0 and G0, which orbit each other every 0.7473 days.
Close to Gamma Corvi and visible in the same binocular field is R Corvi, a long period (Mira) variable star. It ranges in brightness from a magnitude of 6.7 to 14.4 with a period of approximately 317 days. TT Corvi is a semiregular variable red giant of spectral type M3III and apparent magnitude 6.48 around 923 light years distant. It is around 993 times as luminous as the Sun. TU Corvi is a Delta Scuti variable—a class of short period (six hours at most) pulsating stars that have been used as standard candles and as subjects to study astroseismology. It varies by 0.025 of a magnitude around apparent magnitude 6.53 over 59 minutes.
Three star systems have confirmed planets. HD 103774 is a young yellow-white main-sequence star of apparent magnitude 7.12 that is 181 ± 5 light-years distant from Earth. It is 1.335 ± 0.03 times as massive and 3.5 ± 0.3 as luminous as the Sun. Variations in its radial velocity showed it was being orbited by a Neptune-sized planet every 5.9 days in 2013. HD 104067 is an orange dwarf of spectral type K2V of apparent magnitude 7.93 that is 69 ± 1 light-years distant from Earth. Around 80% as massive as the Sun, it is orbited by a planet 3.6 times the mass of Neptune every 55.8 days. WASP-83 has a planet around as massive as Saturn that orbits it every 5 days. It was discovered by its transit across the star in 2015. A fourth star system has an unconfirmed planet. HD 111031 is a sunlike star of spectral type G5V located 101 ± 2 light-years distant from Earth.
Ross 695 is a red dwarf star located a mere 28.9 ± 0.6 light-years distant from Earth. At apparent magnitude 11.27, it is much too faint to be seen with the unaided eye. A small star, it has around 23% the mass and radius of the Sun, but only 0.7% its luminosity. VHS J1256–1257 is a triple system of young brown dwarfs located light-years distant from Earth. The system consists of a central, equal-mass binary system of late-M spectral type dwarfs and an outer, planetary-mass brown dwarf companion that is widely separated at 102 ± 9 AU. DENIS-P J1228.2-1547 is a system composed of two brown dwarfs orbiting each other located 73 ± 3 light-years away from Earth. TV Corvi is a dwarf nova composed of a white dwarf and brown dwarf that orbit each other every 90 minutes. The system has a baseline magnitude of 17 that brightens periodically to magnitude 12, discovered by Clyde Tombaugh in 1931 and David Levy in 1990 and 2005.
Deep-sky objects
Corvus contains no Messier objects. It has several galaxies and a planetary nebula observable with amateur telescopes. The center of Corvus is home to a planetary nebula, NGC 4361. The nebula itself resembles a small elliptical galaxy and has a magnitude of 10.3, but the magnitude 13 star at its centre gives away its true nature. Corvus also contains the Stargate (asterism).
The NGC 4038 Group is a group of galaxies across Corvus and Crater. The group may contain between 13 and 27 galaxies. The best-known member is the Antennae peculiar galaxy, located 0.25 north of 31 Crateris. It consists of two interacting galaxies—NGC 4038 and 4039—that appear to have a heart shape as seen from Earth. The name originates from the huge tidal tails that come off the ends of the two galaxies, formed because of the spiral galaxies' original rotation. Both original galaxies were spiral galaxies and are now experiencing extensive star formation due to the interaction of gas clouds. The galaxies are 45 million light-years from Earth and each has multiple ultraluminous X-ray sources, the source of which is unknown. Astronomers theorize that they may be a rare type of x-ray emitting binary stars or intermediate-mass black holes. The Antennae Galaxies appear in a telescope at the 10th magnitude. SN 2004gt was a type Ic supernova that erupted on December 12, 2004. The progenitor was not identified from older images of the galaxy, and is either a type WC Wolf–Rayet star with a mass over 40 times that of the Sun, or a star 20 to 40 times as massive as the Sun in a binary star system. SN 2007sr was a Type Ia supernova event that peaked in brightness on December 14, 2007. The galaxy has been identified as a good place to take detailed images in case of further supernovae.
NGC 4027 is another member of the NGC 4038 group, notable for its extended spiral arm. Known as the Ringtail Galaxy, it lies close to 31 Crateris. A barred spiral galaxy, its distorted shape is probably due to a past collision, possibly with the nearby NGC 4027A. NGC 4782 and NGC 4783 are a pair of merging elliptical galaxies in the northeastern part of the constellation, around 200 million light-years distant.
Meteor showers
Two established meteor showers originate from within Corvus' boundaries. German astronomer Cuno Hoffmeister discovered and named the Corvids in 1937, after observing them between June 25 and July 2. They have not been seen since, nor was there evidence of a shower when previous records were examined. Hoffmeister noted the trajectory of the shower was similar to that of the comet 11P/Tempel–Swift–LINEAR, though this was not confirmed by Zhukov and colleagues in 2011. The shower has been tentatively linked with 4015 Wilson–Harrington. In January 2013, the MO Video Meteor Network published the discovery of the Eta Corvids, assigning some 300 meteors seen between January 20 and 26. Their existence was confirmed by data analysis later that year.
Popular culture
In 1624, German astronomer Jakob Bartsch equated the constellation Argo Navis with Noah's Ark, linking Corvus and Columba to the crow and dove that feature in the story in Genesis.
In Action Comics #14 (January 2013), which was published 7 November 2012, astrophysicist Neil deGrasse Tyson appears in the story, in which he determines that Superman's home planet, Krypton, orbited the red dwarf LHS 2520 in the constellation Corvus, 27.1 light-years from Earth. Tyson assisted DC Comics in selecting a real-life star that would be an appropriate parent star to Krypton, and picked the star in Corvus, the mascot of Superman's high school, the Smallville Crows.
| Physical sciences | Other | Astronomy |
213468 | https://en.wikipedia.org/wiki/Hydra%20%28constellation%29 | Hydra (constellation) | Hydra is the largest of the 88 modern constellations, measuring 1303 square degrees, and also the longest at over 100 degrees. Its southern end borders Libra and Centaurus and its northern end borders Cancer. It was included among the 48 constellations listed by the 2nd century astronomer Ptolemy. Commonly represented as a water snake, it straddles the celestial equator.
History and mythology
Western mythology
The Greek constellation of Hydra is an adaptation of a Babylonian constellation: the MUL.APIN includes a "serpent" constellation (MUL.DINGIR.MUŠ) that loosely corresponds to Hydra. It is one of two Babylonian "serpent" constellations (the other being the origin of the Greek Serpens), a mythological hybrid of serpent, lion and bird.
The shape of Hydra resembles a twisting snake, and features as such in some Greek myths. One myth associates it with a water snake that a crow served Apollo in a cup when it was sent to fetch water. Apollo saw through the fraud and angrily cast the crow, cup and snake into the sky. It is also associated with the monster Hydra, with its many heads, killed by Hercules, represented in another constellation. According to legend, if one of the Hydra's heads was cut off, two more would grow in its place. However, Hercules' nephew, Iolaus, seared the necks with a torch to prevent them from growing back and thus enabled Hercules to overcome the Hydra.
Mythology and equivalents in other cultures
In Chinese astronomy, the stars that correspond to Hydra are located within the Vermilion Bird and the Azure Dragon. The head of Hydra was collectively known as "Min al Az'al," meaning "belonging to the uninhabited spot" in Arabic.
Features
Stars
Despite its size, Hydra contains only one moderately bright star, Alphard, designated Alpha Hydrae. It is an orange giant of magnitude 2.0, 177 light-years from Earth. Its traditional name means "the solitary one". Beta Hydrae is a blue-white star of magnitude 4.3, 365 light-years from Earth. Gamma Hydrae is a yellow giant of magnitude 3.0, 132 light-years from Earth.
Hydra has one bright binary star, Epsilon Hydrae, which is difficult to split in amateur telescopes; it has a period of 1000 years and is 135 light-years from Earth. The primary is a yellow star of magnitude 3.4 and the secondary is a blue star of magnitude 6.7. However, there are several dimmer double stars and binary stars in Hydra. 27 Hydrae is a triple star with two components visible in binoculars and three visible in small amateur telescopes. The primary is a white star of magnitude 4.8, 244 light-years from Earth. The secondary, a binary star, appears in binoculars at magnitude 7.0 but is composed of a magnitude 7 and a magnitude 11 star; it is 202 light-years from Earth. 54 Hydrae is a binary star 99 light-years from Earth, easily divisible in small amateur telescopes. The primary is a yellow star of magnitude 5.3 and the secondary is a purple star of magnitude 7.4. N Hydrae (N Hya) is a pair of stars of magnitudes 5.8 and 5.9. Struve 1270 (Σ1270) consists of a pair of stars, magnitudes 6.4 and 7.4.
The other main named star in Hydra is Sigma Hydrae (σ Hydrae), which also has the name of Minchir, from the Arabic for snake's nose. At magnitude 4.54, it is rather dim. The head of the snake corresponds to the Āshleshā Nakshatra, the lunar zodiacal constellation in Indian astronomy. The name of Nakshatra (Ashlesha) became the proper name of Epsilon Hydrae since 1 June 2018 by IAU.
Hydra is also home to several variable stars. R Hydrae is a Mira variable star 2000 light-years from Earth; it is one of the brightest Mira variables at its maximum of magnitude 3.5. It has a minimum magnitude of 10 and a period of 390 days. V Hydrae is an unusually vivid red variable star 20,000 light-years from Earth. It varies in magnitude from a minimum of 9.0 to a maximum of 6.6. U Hydrae is a semi-regular variable star with a deep red color, 528 light-years from Earth. It has a minimum magnitude of 6.6 and a maximum magnitude of 4.2; its period is 115 days.
Hydra includes GJ 357, an M-type main sequence star located only 31 light-years from the Solar System. This star has three confirmed exoplanets in its orbit, one of which, GJ 357 d, is considered to be a "Super-Earth" within the circumstellar habitable zone.
The constellation also contains the radio source Hydra A galaxy as well as nearby WISE 0855−0714 brown dwarf being the closest (sub)stellar object of the constellation.
Deep-sky objects
Hydra contains three Messier objects. M83, also known as the Southern Pinwheel Galaxy, is located on the border of Hydra and Centaurus, M68 is a globular cluster near M83, and M48 is an open star cluster in the western end of the serpent.
NGC 3242 is a planetary nebula of magnitude 7.5, 1400 light-years from Earth. Discovered in 1785 by William Herschel, it has earned the nickname "Ghost of Jupiter" because of its striking resemblance to the giant planet. Its blue-green disk is visible in small telescopes and its halo is visible in larger instruments.
M48 (NGC 2548) is an open cluster that is visible to the naked eye under dark skies. Its shape has been described as "triangular"; this 80-star cluster is unusually large, more than half a degree in diameter, larger than the diameter of the full Moon.
There are several globular clusters in Hydra. M68 (NGC 4590) is a globular cluster visible in binoculars and resolvable in medium amateur telescopes. It is 31,000 light-years from Earth and of the 8th magnitude. NGC 5694 is a globular cluster of magnitude 10.2, 105,000 light-years from Earth. Also called "Tombaugh's Globular Cluster", it is a Shapley class VII cluster; the classification indicates that it has intermediate concentration at its nucleus. Though it was discovered as a non-stellar object in 1784 by William Herschel, its status as a globular cluster was not ascertained until 1932, when Clyde Tombaugh looked at photographic plates taken of the region near Pi Hydrae on 12 May 1931.
M83 (NGC 5236), the Southern Pinwheel Galaxy, is an 8th magnitude face-on spiral galaxy. It is easily observed in skies south of 40°N latitude, found by using 1, 2, 3, and 4 Centauri as guide stars. It has been host to six supernovae. Large amateur telescopes - above 12 inches aperture - reveal its spiral arms, bar, and small, bright nucleus. In a medium-sized amateur instrument, around 8 inches in aperture, the spiral arms become visible under good conditions. It is not perfectly symmetrical in the eyepiece, rather, the northwest side is flattened and the nucleus has a southwest-to-northeast bar. A smaller sister to the Milky Way, it is a grand design spiral galaxy 40,000 light-years across.
There are many other galaxies located in Hydra. NGC 3314, usually delineated as NGC 3314a and NGC 3314b, is a pair of galaxies that appear superimposed, despite the fact that they are not related or interacting in any way. The foreground galaxy, NGC 3314a, is at a distance of 140 million light-years, and is a face-on spiral galaxy. The background galaxy, NGC 3314b, is an oblique spiral galaxy, and has a nucleus that appears reddened because of NGC 3314a's dusty disk. ESO 510-G13 is a warped spiral galaxy located 150 million light-years from Earth. Though most galactic disks are flat because of their rate of rotation, their conformation can be changed, as is the case with this galaxy. Astronomers speculate that this is due to interactions with other galaxies. NGC 5068 may be a member of the M83 group, but its identity is disputed. It has a low surface brightness and has a diameter of 4.5 arcminutes. It is not perfectly circular, rather, it is elliptical and oriented on a west-northwest/east-southeast axis. However, it is of fairly uniform brightness throughout.
Another notable galaxy is NGC 4993, an elliptical galaxy which was the source of events GW170817, GRB 170817A and SSS17a from the merger of two neutron stars.
Meteor showers
The Sigma Hydrids peak on December 6 and are a very active shower with an unknown parent body. The Alpha Hydrids are a minor shower that peaks between January 1 and 7.
| Physical sciences | Other | Astronomy |
7750391 | https://en.wikipedia.org/wiki/Collision%20frequency | Collision frequency | Collision frequency describes the rate of collisions between two atomic or molecular species in a given volume, per unit time. In an ideal gas, assuming that the species behave like hard spheres, the collision frequency between entities of species A and species B is:
which has units of [volume][time]−1.
Here,
is the number of A molecules in the gas,
is the number of B molecules in the gas,
is the collision cross section, the "effective area" seen by two colliding molecules, simplified to , where the radius of A and the radius of B.
is the Boltzmann constant,
is the temperature,
is the reduced mass of the reactants A and B,
Collision in diluted solution
In the case of equal-size particles at a concentration in a solution of viscosity , an expression for collision frequency where is the volume in question, and is the number of collisions per second, can be written as:
Where:
is the Boltzmann constant
is the absolute temperature (unit K)
is the viscosity of the solution (pascal seconds)
is the concentration of particles per cm3
Here the frequency is independent of particle size, a result noted as counter-intuitive. For particles of different size, more elaborate expressions can be derived for estimating .
| Physical sciences | Kinetics | Chemistry |
7753430 | https://en.wikipedia.org/wiki/Psychopathy | Psychopathy | Psychopathy, or psychopathic personality, is a personality construct characterized by impaired empathy and remorse, in combination with traits of boldness, disinhibition, and egocentrism. These traits are often masked by superficial charm and immunity to stress, which create an outward appearance of apparent normalcy.
Hervey M. Cleckley, an American psychiatrist, influenced the initial diagnostic criteria for antisocial personality reaction/disturbance in the Diagnostic and Statistical Manual of Mental Disorders (DSM), as did American psychologist George E. Partridge. The DSM and International Classification of Diseases (ICD) subsequently introduced the diagnoses of antisocial personality disorder (ASPD) and dissocial personality disorder (DPD) respectively, stating that these diagnoses have been referred to (or include what is referred to) as psychopathy or sociopathy. The creation of ASPD and DPD was driven by the fact that many of the classic traits of psychopathy were impossible to measure objectively. Canadian psychologist Robert D. Hare later re-popularized the construct of psychopathy in criminology with his Psychopathy Checklist.
Although no psychiatric or psychological organization has sanctioned a diagnosis titled "psychopathy", assessments of psychopathic characteristics are widely used in criminal justice settings in some nations and may have important consequences for individuals. The study of psychopathy is an active field of research. The term is also used by the general public, popular press, and in fictional portrayals. While the abbreviated term "psycho" is often employed in common usage in general media along with "crazy", "insane", and "mentally ill", there is a categorical difference between psychosis and psychopathy.
Signs and symptoms
Socially, psychopathy typically involves extensive callous and manipulative self-serving behaviors with no regard for others and often is associated with repeated delinquency, crime, and violence. Mentally, impairments in processes related to affect and cognition, particularly socially related mental processes, have also been found. Developmentally, symptoms of psychopathy have been identified in young children with conduct disorder and suggest at least a partial constitutional factor that influences its development.
Primary features
Disagreement exists over which features should be considered as part of psychopathy, with researchers identifying around 40 traits supposedly indicative of the construct, though the following characteristics are almost universally considered central.
Core traits
Cooke and Michie (2001) proposed a three-factor model of the Psychopathy Checklist-Revised which has seen widespread application in other measures (e.g., Youth Psychopathic Traits Inventory, Antisocial Process Screening Device).
Arrogant and deceitful interpersonal style: impression management or superficial charm, inflated and grandiose sense of self-worth, pathological lying/deceit, and manipulation for personal gain.
Deficient affective experience: lack of remorse or guilt, shallow affect (coldness and unemotionality), callousness and lack of empathy, and failure to accept responsibility for own actions.
Impulsive and irresponsible lifestyle: impulsivity, sensation-seeking and risk-taking, irresponsible and unreliable behavior, financially parasitic lifestyle, and a lack of realistic, long-term goals.
Low anxiety and fearlessness
Cleckley's (1941) original description of psychopathy included the absence of nervousness and neurotic disorders, and later theorists referred to psychopaths as fearless or thick-skinned. While it is often claimed that the PCL-R does not include low anxiety or fearlessness, such features do contribute to the scoring of the Facet 1 (interpersonal) items, mainly through self-assurance, unrealistic optimism, brazenness, and imperturbability. Indeed, while self-report studies have been inconsistent using the two-factor model of the PCL-R, studies which separate Factor 1 into interpersonal and affective facets, more regularly show modest associations between Facet 1 and low anxiety, boldness and fearless dominance (especially items assessing glibness/charm and grandiosity). When both psychopathy and low anxiety/boldness are measured using interviews, both interpersonal and affective facets are both associated with fearlessness and lack of internalizing disorders.
The importance of low anxiety/fearlessness to psychopathy has historically been underscored through behavioral and physiological studies showing diminished responses to threatening stimuli (interpersonal and affective facets both contributing). However, it is not known whether this is reflected in the reduced experience of state fear or where it reflects impaired detection and response to threat-related stimuli. Moreover, such deficits in threat responding are known to be reduced or even abolished when attention is focused on the threatening stimuli.
Offending
Criminality
In terms of simple correlations, the PCL-R manual states an average score of 22.1 has been found in North American prisoner samples, and that 20.5% scored 30 or higher. An analysis of prisoner samples from outside North America found a somewhat lower average value of 17.5. Studies have found that psychopathy scores correlated with repeated imprisonment, detention in higher security, disciplinary infractions, and substance misuse.
Psychopathy, as measured with the PCL-R in institutional settings, shows in meta-analyses small to moderate effect sizes with institutional misbehavior, postrelease crime, or postrelease violent crime with similar effects for the three outcomes. Individual studies give similar results for adult offenders, forensic psychiatric samples, community samples, and youth. The PCL-R is poorer at predicting sexual re-offending. This small to moderate effect appears to be due largely to the scale items that assess impulsive behaviors and past criminal history, which are well-established but very general risk factors. The aspects of core personality often held to be distinctively psychopathic generally show little or no predictive link to crime by themselves. For example, Factor 1 of the PCL-R and Fearless dominance of the PPI-R have a smaller or no relationship to crime, including violent crime. In contrast, Factor 2 and Impulsive antisociality of the PPI-R are associated more strongly with criminality. Factor 2 has a relationship of similar strength to that of the PCL-R as a whole. The antisocial facet of the PCL-R is still predictive of future violence after controlling for past criminal behavior which, together with results regarding the PPI-R which by design does not include past criminal behavior, suggests that impulsive behaviors is an independent risk factor. Thus, the concept of psychopathy may perform poorly when attempted to be used as a general theory of crime.
Violence
Studies have suggested a strong correlation between psychopathy scores and violence, and the PCL-R emphasizes features that are somewhat predictive of violent behavior. Researchers, however, have noted that psychopathy is dissociable from and not synonymous with violence.
It has been suggested that psychopathy is associated with "instrumental aggression", also known as predatory, proactive, or "cold-blooded" aggression, a form of aggression characterized by reduced emotion and conducted with a goal differing from but facilitated by the commission of harm. One conclusion in this regard was made by a 2002 study of homicide offenders, which reported that the homicides committed by homicidal offenders with psychopathy were almost always (93.3%) primarily instrumental, significantly more than the proportion (48.4%) of those committed by non-psychopathic homicidal offenders, with the instrumentality of the homicide also correlated with the total PCL-R score of the offender as well as their scores on the Factor 1 "interpersonal-affective" dimension. However, contrary to the equating of this to mean exclusively "in cold blood", more than a third of the homicides committed by psychopathic offenders involved some component of emotional reactivity as well. In any case, FBI profilers indicate that serious victim injury is generally an emotional offense, and some research supports this, at least concerning sexual offending. One study has found more serious offending by non-psychopathic offenders on average than by offenders with psychopathy (e.g. more homicides versus more armed robbery and property offenses) and another that the Affective facet of the PCL-R predicted reduced offense seriousness.
Studies on perpetrators of domestic violence find that abusers have high rates of psychopathy, with the prevalence estimated to be at around 15-30%. Furthermore, the commission of domestic violence is correlated with Factor 1 of the PCL-R, which describes the emotional deficits and the callous and exploitative interpersonal style found in psychopathy. The prevalence of psychopathy among domestic abusers indicate that the core characteristics of psychopathy, such as callousness, remorselessness, and a lack of close interpersonal bonds, predispose those with psychopathy to committing domestic abuse, and suggest that the domestic abuses committed by these individuals are callously perpetrated (i.e. instrumentally aggressive) rather than a case of emotional aggression and therefore may not be amenable to the types of psychosocial interventions commonly given to domestic abuse perpetrators.
Some clinicians suggest that assessment of the construct of psychopathy does not necessarily add value to violence risk assessment. A large systematic review and meta-regression found that the PCL performed the poorest out of nine tools for predicting violence. In addition, studies conducted by the authors or translators of violence prediction measures, including the PCL, show on average more positive results than those conducted by more independent investigators. Several other risk assessment instruments can predict further crime with an accuracy similar to the PCL-R and some of these are considerably easier, quicker, and less expensive to administer. This may even be done automatically by a computer simply based on data such as age, gender, number of previous convictions, and age of first conviction. Some of these assessments may also identify treatment changes and goals, identify quick changes that may help short-term management, identify more specific kinds of violence that may be at risk, and establish specific probabilities of offending for specific scores. Nonetheless, the PCL-R may continue to be popular for risk assessment because of its pioneering role and the large amount of research done using it.
The U.S. Federal Bureau of Investigation (FBI) reports that psychopathic behavior is consistent with traits common to some serial killers, including sensation seeking, a lack of remorse or guilt, impulsivity, the need for control, and predatory behavior. It has also been found that the homicide victims of psychopathic offenders were disproportionately female in comparison to the more equitable gender distribution of victims of non-psychopathic offenders.
Sexual offending
Psychopathy has been associated with the commission of sexual crime, with some researchers arguing that it is correlated with a preference for violent sexual behavior. A 2011 study of conditional releases for Canadian male federal offenders found that psychopathy was related to more violent and non-violent offences but not more sexual offences. For child molesters, psychopathy was associated with more offences. A study on the relationship between psychopathy scores and types of aggression in a sample of sexual murderers, in which 84.2% of the sample had PCL-R scores above 20 and 47.4% above 30, found that 82.4% of those with scores above 30 had engaged in sadistic violence (defined as enjoyment indicated by self-report or evidence) compared to 52.6% of those with scores below 30, and total PCL-R and Factor 1 scores correlated significantly with sadistic violence. Despite this, it is reported that offenders with psychopathy (both sexual and non-sexual offenders) are about 2.5 times more likely to be granted conditional release compared to non-psychopathic offenders.
Hildebrand and colleagues (2004) have uncovered an interaction between psychopathy and deviant sexual interests, wherein those high in psychopathy who also endorsed deviant sexual interests were more likely to recidivate sexually. A subsequent meta-analysis has consolidated such a result.
In considering the issue of possible reunification of some sex offenders into homes with a non-offending parent and children, it has been advised that any sex offender with a significant criminal history should be assessed on the PCL-R, and if they score 18 or higher, then they should be excluded from any consideration of being placed in a home with children under any circumstances. There is, however, increasing concern that PCL scores are too inconsistent between different examiners, including in its use to evaluate sex offenders.
Other offending
The possibility of psychopathy has been associated with organized crime, economic crime, and war crimes. Terrorists are sometimes considered psychopathic, and comparisons may be drawn with traits such as antisocial violence, a selfish worldview that precludes the welfare of others, a lack of remorse or guilt, and blame externalization. However, John Horgan, author of The Psychology of Terrorism, argues that such comparisons could also then be drawn more widely: for example, to soldiers in wars. Coordinated terrorist activity requires organization, loyalty, and ideological fanaticism often to the extreme of sacrificing oneself for an ideological cause. Traits such as a self-centered disposition, unreliability, poor behavioral controls, and unusual behaviors may disadvantage or preclude psychopathic individuals from conducting organized terrorism.
It may be that a significant portion of people with psychopathy are socially successful and tend to express their antisocial behavior through more covert avenues such as social manipulation or white collar crime. Such individuals are sometimes referred to as "successful psychopaths", and may not necessarily always have extensive histories of traditional antisocial behavior as characteristic of traditional psychopathy.
Childhood and adolescent precursors
The PCL:YV is an adaptation of the PCL-R for individuals aged 13–18 years. It is, like the PCL-R, done by a trained rater based on an interview and an examination of criminal and other records. The "Antisocial Process Screening Device" (APSD) is also an adaptation of the PCL-R. It can be administered by parents or teachers for individuals aged 6–13 years. High psychopathy scores for both juveniles (as measured with these instruments) and adults (as measured with the PCL-R and other measurement tools) have similar associations with other variables, including similar ability in predicting violence and criminality. Juvenile psychopathy may also be associated with more negative emotionality such as anger, hostility, anxiety, and depression. Psychopathic traits in youth typically comprise three factors: callous/unemotional, narcissism, and impulsivity/irresponsibility.
There is a positive correlation between early negative life events of the ages 0–4 and the emotion-based aspects of psychopathy. There are moderate to high correlations between psychopathy rankings from late childhood to early adolescence. The correlations are considerably lower from early- or mid-adolescence to adulthood. In one study most of the similarities were on the Impulsive- and Antisocial-Behavior scales. Of those adolescents who scored in the top 5% highest psychopathy scores at age 13, less than one-third (29%) were classified as psychopathic at age 24. Some recent studies have also found a poorer ability to predict long-term, adult offending.
Conduct disorder
Conduct disorder is diagnosed based on a prolonged pattern of antisocial behavior in childhood and/or adolescence, and may be seen as a precursor to ASPD. Some researchers have speculated that two subtypes of conduct disorder mark dual developmental pathways to adult psychopathy. The DSM allows differentiating between childhood onset before age 10 and adolescent onset at age 10 and later. Childhood-onset is argued to be more due to a personality disorder caused by neurological deficits interacting with an adverse environment. For many, but not all, childhood onset is associated with what is in Terrie Moffitt's developmental theory of crime referred to as "life-course- persistent" antisocial behavior as well as poorer health and economic status. Adolescent onset is argued to more typically be associated with short-term antisocial behavior.
It has been suggested that the combination of early-onset conduct disorder and ADHD may be associated with life-course-persistent antisocial behaviors as well as psychopathy. There is evidence that this combination is more aggressive and antisocial than those with conduct disorder alone. However, it is not a particularly distinct group since the vast majority of young children with conduct disorder also have ADHD. Some evidence indicates that this group has deficits in behavioral inhibition, similar to that of adults with psychopathy. They may not be more likely than those with conduct disorder alone to have the interpersonal/affective features and the deficits in emotional processing characteristic of adults with psychopathy. Proponents of different types/dimensions of psychopathy have seen this type as possibly corresponding to adult secondary psychopathy and increased disinhibition in the triarchic model.
The DSM-5 includes a specifier for those with conduct disorder who also display a callous, unemotional interpersonal style across multiple settings and relationships. The specifier is based on research that suggests that those with conduct disorder who also meet the criteria for the specifier tend to have a more severe form of the disorder with an earlier onset as well as a different response to treatment. Proponents of different types/dimensions of psychopathy have seen this as possibly corresponding to adult primary psychopathy and increased boldness and/or meanness in the triarchic model.
Mental traits
Cognition
Dysfunctions in the prefrontal cortex and amygdala regions of the brain have been associated with specific learning impairments in psychopathy. Damage to the ventromedial prefrontal cortex, which regulates the activity in the amygdala, leads to common characteristics in psychopathic individuals. Since the 1980s, scientists have linked traumatic brain injury, including damage to these regions, with violent and psychopathic behavior. Patients with damage in such areas resembled "psychopathic individuals" whose brains were incapable of acquiring social and moral knowledge; those who acquired damage as children may have trouble conceptualizing social or moral reasoning, while those with adult-acquired damage may be aware of proper social and moral conduct but be unable to behave appropriately. Dysfunctions in the amygdala and ventromedial prefrontal cortex may also impair stimulus-reinforced learning in psychopaths, whether punishment-based or reward-based. People scoring 25 or higher in the PCL-R, with an associated history of violent behavior, appear to have significantly reduced mean microstructural integrity in their uncinate fasciculus—white matter connecting the amygdala and orbitofrontal cortex. There is evidence from DT-MRI of breakdowns in the white matter connections between these two important areas.
Although some studies have suggested inverse relationships between psychopathy and intelligence, including with regards to verbal IQ, Hare and Neumann state that a large literature demonstrates at most only a weak association between psychopathy and IQ, noting that the early pioneer Cleckley included good intelligence in his checklist due to selection bias (since many of his patients were "well educated and from middle-class or upper-class backgrounds") and that "there is no obvious theoretical reason why the disorder described by Cleckley or other clinicians should be related to intelligence; some psychopaths are bright, others less so". Studies also indicate that different aspects of the definition of psychopathy (e.g. interpersonal, affective (emotion), behavioral, and lifestyle components) can show different links to intelligence, and the result can depend on the type of intelligence assessment (e.g. verbal, creative, practical, analytical).
Emotion recognition and empathy
A large body of research suggests that psychopathy is associated with atypical responses to distress cues from other people, more precisely an impaired emotional empathy in the recognition of, and response to, facial expressions, body gestures and vocal tones of fear, sadness, pain and happiness. This impaired recognition and reduced autonomic responsiveness might be partly accounted for by a decreased activation of the fusiform and extrastriate cortical regions. The underlying biological surfaces for processing expressions of happiness are functionally intact in psychopaths, although less responsive than those of controls. The neuroimaging literature is unclear as to whether deficits are specific to particular emotions such as fear. The overall pattern of results across studies indicates that people diagnosed with psychopathy demonstrate reduced MRI, fMRI, aMRI, PET, and SPECT activity in areas of the brain. Research has also shown that an approximate 18% smaller amygdala size contributes to a significantly lower emotional sensation in regards to fear, sadness, amongst other negative emotions, which may likely be the reason as to why psychopathic individuals have lower empathy. Some recent fMRI studies have reported that emotion perception deficits in psychopathy are pervasive across emotions (positives and negatives). Studies on children with psychopathic tendencies have also shown such associations. Meta-analyses have also found evidence of impairments in both vocal and facial emotional recognition for several emotions (i.e., not only fear and sadness) in both adults and children/adolescents.
Moral judgment
Psychopathy has been associated with amorality—an absence of, indifference towards, or disregard for moral beliefs. There are few firm data on patterns of moral judgment. Studies of the developmental level (sophistication) of moral reasoning found all possible results—lower, higher or the same as non-psychopaths. Studies that compared judgments of personal moral transgressions versus judgments of breaking conventional rules or laws found that psychopaths rated them as equally severe, whereas non-psychopaths rated the rule-breaking as less severe.
A study comparing judgments of whether personal or impersonal harm would be endorsed to achieve the rational maximum (utilitarian) amount of welfare found no significant differences between subjects high and low in psychopathy. However, a further study using the same tests found that prisoners scoring high on the PCL were more likely to endorse impersonal harm or rule violations than non-psychopathic controls were. The psychopathic offenders who scored low in anxiety were also more willing to endorse personal harm on average.
Assessing accidents, where one person harmed another unintentionally, psychopaths judged such actions to be more morally permissible. This result has been considered a reflection of psychopaths' failure to appreciate the emotional aspect of the victim's harmful experience.
History
Etymology
The word psychopathy is a joining of the Greek words psyche () "soul" and pathos () "suffering, feeling". The first documented use is from 1847 in Germany as psychopatisch, and the noun psychopath has been traced to 1885. In medicine, patho- has a more specific meaning of disease (Thus pathology has meant the study of disease since 1610, and psychopathology has meant the study of mental disorder in general since 1847. A sense of "a subject of pathology, morbid, excessive" is attested from 1845, including the phrase pathological liar from 1891 in the medical literature).
The term psychopathy initially had a very general meaning referring to all sorts of mental disorders and social aberrations, popularised from 1891 in Germany by Koch's concept of "psychopathic inferiority" (). Some medical dictionaries still define psychopathy in both a narrow and broad sense, such as MedlinePlus from the U.S. National Library of Medicine. On the other hand, Stedman's Medical Dictionary defines "psychopath" only as a "former designation" for a person with an antisocial type of personality disorder.
The term psychosis was also used in Germany from 1841, originally in a very general sense. The suffix -ωσις (-osis) meant in this case "abnormal condition". This term or its adjective psychotic would come to refer to the more severe mental disturbances and then specifically to mental states or disorders characterized by hallucinations, delusions, or in some other sense markedly out of touch with reality.
The slang term psycho has been traced to a shortening of the adjective psychopathic from 1936, and from 1942 as a shortening of the noun psychopath, but it is also used as shorthand for psychotic or crazed.
The media usually uses the term psychopath to designate any criminal whose offenses are particularly abhorrent and unnatural, but that is not its original or general psychiatric meaning.
Sociopathy
The word element socio- has been commonly used in compound words since around 1880. The term sociopathy may have been first introduced in 1909 in Germany by biological psychiatrist Karl Birnbaum and in 1930 in the US by educational psychologist George E. Partridge, as an alternative to the concept of psychopathy. It was used to indicate that the defining feature is violation of social norms, or antisocial behavior, and may be social or biological in origin.
The terms sociopathy and psychopathy were once used interchangeably concerning antisocial personality disorder, though this usage is outdated in medicine and psychiatry. Psychopathy, however, is a highly popular construct in the psychology literature. Furthermore, the DSM-5 introduced the dimensional model of personality disorders in Section III, which includes a specifier for psychopathic traits. According to the DSM, psychopathy is not a standalone diagnosis, but the authors attempted to measure "psychopathic traits" via a specifier. In one study, the "Psychopathic Features Specifier" has been modeled on Factor 1 of the Psychopathic Personality Inventory, known as Fearless Dominance. To some, it is evidence of psychopathy not being a more extreme version of ASPD, but as an emergent compound trait that manifests when Antisocial Personality Disorder is present in combination with high levels of Fearless Dominance (or Boldness as it's known in the Triarchic Model). Analyses showed that this Section III ASPD greatly outperformed Section II ASPD in predicting scores on Hare’s (2003) Psychopathy Checklist-Revised.
Section III ASPD including the 'Psychopathic Traits Specifier' can be seen on page 765 of the DSM-5 or Page 885 of the DSM-5-TR.
The term is used in various ways in contemporary usage. Robert Hare stated in the popular science book Snakes in Suits that sociopathy and psychopathy are often used interchangeably, but in some cases the term sociopathy is preferred because it is less likely than is psychopathy to be confused with psychosis, whereas in other cases the two terms may be used with different meanings that reflect the user's views on its origins and determinants. Hare contended that the term sociopathy is preferred by those who see the causes as due to social factors and early environment, and the term psychopathy is preferred by those who believe that there are psychological, biological, and genetic factors involved in addition to environmental factors. Hare also provides his own definitions: he describes psychopathy as lacking a sense of empathy or morality, but sociopathy as only differing from the average person in the sense of right and wrong.
Precursors
Ancient writings that have been connected to psychopathic traits include Deuteronomy and a description of an unscrupulous man by the Greek philosopher Theophrastus around 300 BC.
The concept of psychopathy has been indirectly connected to the early 19th-century work of Pinel (1801; "mania without delirium") and Pritchard (1835; "moral insanity"), although historians have largely discredited the idea of a direct equivalence. Psychopathy originally described any illness of the mind, but found its application to a narrow subset of mental conditions when it was used toward the end of the 19th century by the German psychiatrist Julius Koch (1891) to describe various behavioral and moral dysfunction in the absence of an obvious mental illness or intellectual disability. He applied the term psychopathic inferiority () to various chronic conditions and character disorders, and his work would influence the later conception of the personality disorder.
The term psychopathic came to be used to describe a diverse range of dysfunctional or antisocial behavior and mental and sexual deviances, including at the time homosexuality. It was often used to imply an underlying "constitutional" or genetic origin. Disparate early descriptions likely set the stage for modern controversies about the definition of psychopathy.
20th century
An influential figure in shaping modern American conceptualizations of psychopathy was American psychiatrist Hervey Cleckley. In his classic monograph, The Mask of Sanity (1941), Cleckley drew on a small series of vivid case studies of psychiatric patients at a Veterans Administration hospital in Georgia to describe psychopathy. Cleckley used the metaphor of the "mask" to refer to the tendency of psychopaths to appear confident, personable, and well-adjusted compared to most psychiatric patients while revealing underlying pathology through their actions over time. Cleckley formulated sixteen criteria for psychopathy. The Scottish psychiatrist David Henderson had also been influential in Europe from 1939 in narrowing the diagnosis.
The diagnostic category of sociopathic personality in early editions of the Diagnostic and Statistical Manual (DSM) had some key similarities to Cleckley's ideas, though in 1980 when renamed Antisocial Personality Disorder some of the underlying personality assumptions were removed. In 1980, Canadian psychologist Robert D. Hare introduced an alternative measure, the "Psychopathy Checklist" (PCL) based largely on Cleckley's criteria, which was revised in 1991 (PCL-R), and is the most widely used measure of psychopathy. There are also several self-report tests, with the Psychopathic Personality Inventory (PPI) used more often among these in contemporary adult research.
Famous individuals have sometimes been diagnosed, albeit at a distance, as psychopaths. As one example out of many possible from history, in a 1972 version of a secret report originally prepared for the Office of Strategic Services in 1943, which may have been intended to be used as propaganda, non-medical psychoanalyst Walter C. Langer suggested Adolf Hitler was probably a psychopath. However, others have not drawn this conclusion; clinical forensic psychologist Glenn Walters argues that Hitler's actions do not warrant a diagnosis of psychopathy as, although he showed several characteristics of criminality, he was not always egocentric, callously disregarding of feelings or lacking impulse control, and there is no proof he could not learn from mistakes.
Definition
Concepts
There are multiple conceptualizations of psychopathy, including Cleckleyan psychopathy (Hervey Cleckley's conception entailing bold, disinhibited behavior, and "feckless disregard") and criminal psychopathy (a meaner, more aggressive and disinhibited conception explicitly entailing persistent and sometimes serious criminal behavior). The latter conceptualization is typically used as the modern clinical concept and assessed by the Psychopathy Checklist. The label "psychopath" may have implications and stigma related to decisions about punishment severity for criminal acts, medical treatment, civil commitments, etc. Efforts have therefore been made to clarify the meaning of the term.
It has been suggested that those who share the same emotional deficiencies and psychopathic features, but are properly socialized, should not be designated as 'psychopaths'.
The triarchic model suggests that different conceptions of psychopathy emphasize three observable characteristics to various degrees. Analyses have been made with respect to the applicability of measurement tools such as the Psychopathy Checklist (PCL, PCL-R) and Psychopathic Personality Inventory (PPI) to this model.
Boldness. Low fear including stress-tolerance, toleration of unfamiliarity and danger, and high self-confidence and social assertiveness. The PCL-R measures this relatively poorly and mainly through Facet 1 of Factor 1. Similar to PPI fearless dominance. May correspond to differences in the amygdala and other neurological systems associated with fear.
Disinhibition. Poor impulse control including problems with planning and foresight, lacking affect and urge control, demand for immediate gratification, and poor behavioral restraints. Similar to PCL-R Factor 2 and PPI impulsive antisociality. May correspond to impairments in frontal lobe systems that are involved in such control.
Meanness. Lacking empathy and close attachments with others, disdain of close attachments, use of cruelty to gain empowerment, exploitative tendencies, defiance of authority, and destructive excitement seeking. The PCL-R in general is related to this but in particular some elements in Factor 1. Similar to PPI, but also includes elements of subscales in impulsive antisociality.
Psychopathy has been conceptualized as a hybrid condition marked by a paradoxical combination of superficial charm, poise, emotional resilience, and venturesomeness on the outside but deep-seated affective disturbances and impulse control deficits on the inside. From this perspective, psychopathy is at least in part characterized by psychologically adaptive traits. Furthermore, according to this view, psychopathy may be linked to at least some interpersonally successful outcomes, such as effective leadership, business accomplishments, and heroism.
Measurement
An early and influential analysis from Harris and colleagues indicated that a discrete category, or taxon, may underlie PCL-R psychopathy, allowing it to be measured and analyzed. However, this was only found for the behavioral Factor 2 items they identified, child problem behaviors; adult criminal behavior did not support the existence of a taxon. Marcus, John, and Edens more recently performed a series of statistical analyses on PPI scores and concluded that psychopathy may best be conceptualized as having a "dimensional latent structure" like depression.
Marcus et al. repeated the study on a larger sample of prisoners, using the PCL-R and seeking to rule out other experimental or statistical issues that may have produced the previously different findings. They again found that the psychopathy measurements do not appear to be identifying a discrete type (a taxon). They suggest that while for legal or other practical purposes an arbitrary cut-off point on trait scores might be used, there is actually no clear scientific evidence for an objective point of difference by which to label some people "psychopaths"; in other words, a "psychopath" may be more accurately described as someone who is "relatively psychopathic".
The PCL-R was developed for research, not clinical forensic diagnosis, and even for research purposes to improve understanding of the underlying issues, it is necessary to examine dimensions of personality in general rather than only a constellation of traits. The PCL-R test has been used to determine "true" or primary psychopaths (individuals that score a 30 or higher on the PCL-R test). Primary psychopaths are distinguished from secondary psychopaths, and contrast with those who are legitimately considered antisocial.
Personality dimensions
Studies have linked psychopathy to alternative dimensions such as antagonism (high), conscientiousness (low), and anxiousness (low).
Psychopathy has also been linked to high psychoticism—a theorized dimension referring to tough, aggressive, or hostile tendencies. Aspects of this that appear associated with psychopathy are lack of socialization and responsibility, impulsivity, sensation-seeking (in some cases), and aggression.
Otto Kernberg, from a particular psychoanalytic perspective, believed psychopathy should be considered as part of a spectrum of pathological narcissism, that would range from narcissistic personality on the low end, malignant narcissism in the middle, and psychopathy at the high end.
Psychopathy, narcissism, and Machiavellianism, three personality traits that are together referred to as the dark triad, share certain characteristics, such as a callous-manipulative interpersonal style. The dark tetrad refers to these traits with the addition of sadism. Several psychologists have asserted that subclinical psychopathy and Machiavellianism are more or less interchangeable. There is a subscale on the Psychopathic Personality Inventory (PPI) dubbed "Machiavellian Egocentricity". Delroy Paulhus has asserted that the difference that most miss is that while both are characterized by manipulativeness and unemotionality, psychopaths tend to be more reckless. One study asserted that "the
ability to adapt, reappraise and reassess a situation may be key factors differentiating
Machiavellianism from psychopathy, for example". Psychopathy and Machiavellianism were also correlated similarly in responses to affective stimuli, and both are negatively correlated with the recognition of facial emotions. Many have suggested merging the dark triad traits (especially Machiavellianism and psychopathy) into one construct, given empirical studies that show immense overlap.
Criticism of current conceptions
The current conceptions of psychopathy have been criticized for being poorly conceptualized, highly subjective, and encompassing a wide variety of underlying disorders. Dorothy Otnow Lewis has written:
Half of the Hare Psychopathy Checklist consists of symptoms of mania, hypomania, and frontal-lobe dysfunction, which frequently results in underlying disorders being dismissed. Hare's conception of psychopathy has also been criticized for being reductionist, dismissive, tautological, and ignorant of context as well as the dynamic nature of human behavior. Some have called for rejection of the concept altogether, due to its vague, subjective and judgmental nature that makes it prone to misuse. A systematic review determined that the PCL is weakly predictive of criminal behavior, but not of lack of conscience, or treatment and rehabilitation outcomes. These findings contradict widespread beliefs among professionals in forensics.
Psychopathic individuals do not show regret or remorse. This was thought to be due to an inability to generate this emotion in response to negative outcomes. However, in 2016, people with antisocial personality disorder were found to experience regret but did not use the regret to guide their choice in behavior. There was no lack of regret but a problem to think through a range of potential actions and estimating the outcome values.
In an experiment published in March 2007 at the University of Southern California neuroscientist Antonio R. Damasio and his colleagues showed that subjects with damage to the ventromedial prefrontal cortex lack the ability to empathically feel their way to moral answers, and that when confronted with moral dilemmas, these brain-damaged patients coldly came up with "end-justifies-the-means" answers, leading Damasio to conclude that the point was not that they reached immoral conclusions, but that when they were confronted by a difficult issue – in this case as whether to shoot down a passenger plane hijacked by terrorists before it hits a major city – these patients appear to reach decisions without the anguish that afflicts those with typically functioning brains. According to Adrian Raine, a clinical neuroscientist also at the University of Southern California, one of this study's implications is that society may have to rethink how it judges immoral people: "Psychopaths often feel no empathy or remorse. Without that awareness, people relying exclusively on reasoning seem to find it harder to sort their way through moral thickets. Does that mean they should be held to different standards of accountability?"
Cause
Behavioral genetic studies have identified potential genetic and non-genetic contributors to psychopathy, including influences on brain function. Proponents of the triarchic model believe that psychopathy results from the interaction of genetic predispositions and an adverse environment. What is adverse may differ depending on the underlying predisposition: for example, it is hypothesized that persons having high boldness may respond poorly to punishment but may respond better to rewards and secure attachments.
Genetic
Genetically informed studies of the personality characteristics typical of individuals with psychopathy have found moderate genetic (as well as non-genetic) influences. On the PPI, fearless dominance and impulsive antisociality were similarly influenced by genetic factors and uncorrelated with each other. Genetic factors may generally influence the development of psychopathy while environmental factors affect the specific expression of the traits that predominate. A study on a large group of children found more than 60% heritability for "callous-unemotional traits" and that conduct disorder among children with these traits has a higher heritability than among children without these traits.
Environment
A study by Farrington of a sample of London males between ages 8 and 48 included studying which factors scored 10 or more on the PCL:SV at age 48. The strongest factors included having a convicted parent, being physically neglected, low involvement of the father with the boy, low family income, and coming from a disrupted family. Other significant factors included poor supervision, abuse, harsh discipline, large family size, delinquent siblings, young mothers, depressed mothers, low social class, and poor housing. There has also been an association between psychopathy and detrimental treatment by peers. However, it is difficult to determine the extent of an environmental influence on the development of psychopathy because of evidence of its strong heritability.
Brain injury
Researchers have linked head injuries with psychopathy and violence. Since the 1980s, scientists have associated traumatic brain injury, such as damage to the prefrontal cortex, including the orbitofrontal cortex, with psychopathic behavior and a deficient ability to make morally and socially acceptable decisions, a condition that has been termed "acquired sociopathy", or "pseudopsychopathy". Individuals with damage to the area of the prefrontal cortex known as the ventromedial prefrontal cortex show remarkable similarities to diagnosed psychopathic individuals, displaying reduced autonomic response to emotional stimuli, deficits in aversive conditioning, similar preferences in moral and economic decision making, and diminished empathy and social emotions like guilt or shame. These emotional and moral impairments may be especially severe when the brain injury occurs at a young age. Children with early damage in the prefrontal cortex may never fully develop social or moral reasoning and become "psychopathic individuals ... characterized by high levels of aggression and antisocial behavior performed without guilt or empathy for their victims". Additionally, damage to the amygdala may impair the ability of the prefrontal cortex to interpret feedback from the limbic system, which could result in uninhibited signals that manifest in violent and aggressive behavior.
Childhood trauma
Other theories
Evolutionary explanations
Psychopathy is associated with several adverse life outcomes as well as an increased risk of disability and death due to factors such as violence, accidents, homicides, and suicides. This, in combination with the evidence for genetic influences, is evolutionarily puzzling and may suggest that there are compensating evolutionary advantages, and researchers within evolutionary psychology have proposed several evolutionary explanations. According to one hypothesis, some traits associated with psychopathy may be socially adaptive, and psychopathy may be a frequency-dependent, socially parasitic strategy, which may work as long as there is a large population of altruistic and trusting individuals, relative to the population of psychopathic individuals, to be exploited. It is also suggested that some traits associated with psychopathy such as early, promiscuous, adulterous, and coercive sexuality may increase reproductive success. Robert Hare has stated that many psychopathic males have a pattern of mating with and quickly abandoning women, and thereby have a high fertility rate, resulting in children that may inherit a predisposition to psychopathy.
Criticism includes that it may be better to look at the contributing personality factors rather than treat psychopathy as a unitary concept due to poor testability. Furthermore, if psychopathy is caused by the combined effects of a very large number of adverse mutations then each mutation may have such a small effect that it escapes natural selection. The personality is thought to be influenced by a very large number of genes and may be disrupted by random mutations, and psychopathy may instead be a product of a high mutation load. Psychopathy has alternatively been suggested to be a spandrel, a byproduct, or side-effect, of the evolution of adaptive traits rather than an adaptation in itself.
Mechanisms
Psychological
Some laboratory research demonstrates correlations between psychopathy and atypical responses to aversive stimuli, including weak conditioning to painful stimuli and poor learning of avoiding responses that cause punishment, as well as low reactivity in the autonomic nervous system as measured with skin conductance while waiting for a painful stimulus but not when the stimulus occurs. While it has been argued that the reward system functions normally, some studies have also found reduced reactivity to pleasurable stimuli. According to the response modulation hypothesis, psychopathic individuals have also had difficulty switching from an ongoing action despite environmental cues signaling a need to do so. This may explain the difficulty responding to punishment, although it is unclear if it can explain findings such as deficient conditioning. There may be methodological issues regarding the research. While establishing a range of idiosyncrasies on average in linguistic and affective processing under certain conditions, this research program has not confirmed a common pathology of psychopathy.
Neurological
Thanks to advancing MRI studies, experts can visualize specific brain differences and abnormalities of individuals with psychopathy in areas that control emotions, social interactions, ethics, morality, regret, impulsivity, and conscience within the brain. Blair, a researcher who pioneered research into psychopathic tendencies stated, "With regard to psychopathy, we have clear indications regarding why the pathology gives rise to the emotional and behavioral disturbance and important insights into the neural systems implicated in this pathology". Dadds et al., remarks that despite a rapidly advancing neuroscience of empathy, little is known about the developmental underpinnings of the psychopathic disconnect between affective and cognitive empathy.
A 2008 review by Weber et al. suggested that psychopathy is sometimes associated with brain abnormalities in prefrontal-temporo-limbic regions that are involved in emotional and learning processes, among others. Neuroimaging studies have found structural and functional differences between those scoring high and low on the PCL-R in a 2011 review by Skeem et al. stating that they are "most notably in the amygdala, hippocampus and parahippocampal gyri, anterior and posterior cingulate cortex, striatum, insula, and frontal and temporal cortex".
The amygdala and frontal areas have been suggested as particularly important. People scoring 25 or higher in the PCL-R, with an associated history of violent behavior, appear on average to have significantly reduced microstructural integrity between the white matter connecting the amygdala and orbitofrontal cortex (such as the uncinate fasciculus). The evidence suggested that the degree of abnormality was significantly related to the degree of psychopathy and may explain the offending behaviors. Furthermore, changes in the amygdala have been associated with "callous-unemotional" traits in children. However, the amygdala has also been associated with positive emotions, and there have been inconsistent results in the studies in particular areas, which may be due to methodological issues. Others have cast doubt on the amygdala as important for psychopathy, with one meta-analysis suggesting that most studies on the amygdala and psychopathy find no effect and that studies finding a negative effect (that psychopaths display less amygdala activity) have lower statistical power.
Some of these findings are consistent with other research and theories. For example, in a neuroimaging study of how individuals with psychopathy respond to emotional words, widespread differences in activation patterns have been shown across the temporal lobe when psychopathic criminals were compared to "normal" volunteers, which is consistent with views in clinical psychology. Additionally, the notion of psychopathy being characterized by low fear is consistent with findings of abnormalities in the amygdala, since deficits in aversive conditioning and instrumental learning are thought to result from amygdala dysfunction, potentially compounded by orbitofrontal cortex dysfunction, although the specific reasons are unknown.
Considerable research has documented the presence of the two subtypes of primary and secondary psychopathy. Proponents of the primary-secondary psychopathy distinction and triarchic model argue that there are neurological differences between these subgroups of psychopathy which support their views. For instance, the boldness factor in the triarchic model is argued to be associated with reduced activity in the amygdala during fearful or aversive stimuli and reduced startle response, while the disinhibition factor is argued to be associated with impairment of frontal lobe tasks. There is evidence that boldness and disinhibition are genetically distinguishable.
Biochemical
High levels of testosterone combined with low levels of cortisol and/or serotonin have been theorized as contributing factors. Testosterone is "associated with approach-related behavior, reward sensitivity, and fear reduction", and injecting testosterone "shift[s] the balance from punishment to reward sensitivity", decreases fearfulness, and increases "responding to angry faces". Some studies have found that high testosterone levels are associated with antisocial and aggressive behaviors, yet other research suggests that testosterone alone does not cause aggression but increases dominance-seeking. It is unclear from studies if psychopathy correlates with high testosterone levels, but a few studies have found that disruption of serotonin neurotransmission disrupts cortisol reactivity to a stress-inducing speech task. Thus, dysregulation of serotonin in the brain may contribute to the low cortisol levels observed in psychopathy. Cortisol increases withdrawal behavior and sensitivity to punishment and aversive conditioning, which are abnormally low in individuals with psychopathy and may underlie their impaired aversion learning and disinhibited behavior. High testosterone levels combined with low serotonin levels are associated with "impulsive and highly negative reactions", and may increase violent aggression when an individual is provoked or becomes frustrated. Several animal studies note the role of serotonergic functioning in impulsive aggression and antisocial behavior.
However, some studies on animal and human subjects have suggested that the emotional-interpersonal traits and predatory aggression of psychopathy, in contrast to impulsive and reactive aggression, is related to increased serotoninergic functioning. A study by Dolan and Anderson, regarding the relationship between serotonin and psychopathic traits in a sample of personality disordered offenders, found that serotonin functioning as measured by prolactin response, while inversely associated with impulsive and antisocial traits, were positively correlated with arrogant and deceitful traits, and, to a lesser extent, callous and remorseless traits. Bariş Yildirim theorizes that the 5-HTTLPR "long" allele, which is generally regarded as protective against internalizing disorders, may interact with other serotoninergic genes to create a hyper-regulation and dampening of affective processes that results in psychopathy's emotional impairments. Furthermore, the combination of the 5-HTTLPR long allele and high testosterone levels has been found to result in a reduced response to threat as measured by cortisol reactivity, which mirrors the fear deficits found in those with psychopathy.
Studies have suggested other correlations. Psychopathy was associated in two studies with an increased ratio of HVA (a dopamine metabolite) to 5-HIAA (a serotonin metabolite). Studies have found that individuals with the traits meeting criteria for psychopathy show a greater dopamine response to potential "rewards" such as monetary promises or taking drugs such as amphetamines. This has been theoretically linked to increased impulsivity. A 2010 British study found that a large 2D:4D digit ratio, an indication of high prenatal estrogen exposure, was a "positive correlate of psychopathy in females, and a positive correlate of callous affect (psychopathy sub-scale) in males".
Findings have also shown monoamine oxidase A to affect the predictive ability of the PCL-R. Monoamine oxidases (MAOs) are enzymes that are involved in the breakdown of neurotransmitters such as serotonin and dopamine and are, therefore, capable of influencing feelings, mood, and behavior in individuals. Findings suggest that further research is needed in this area.
Diagnosis
Tools
Psychopathy Checklist
Psychopathy is most commonly assessed with the Psychopathy Checklist, Revised (PCL-R), created by Robert D. Hare based on Cleckley's criteria from the 1940s, criminological concepts such as those of William and Joan McCord, and his own research on criminals and incarcerated offenders in Canada. The PCL-R is widely used and is referred to by some as the "gold standard" for assessing psychopathy. There are nonetheless numerous criticisms of the PCL-R as a theoretical tool and in real-world usage.
Psychopathic Personality Inventory
Unlike the PCL, the Psychopathic Personality Inventory (PPI) was developed to comprehensively index personality traits without explicitly referring to antisocial or criminal behaviors themselves. It is a self-report scale that was developed originally for non-clinical samples (e.g. university students) rather than prisoners, though may be used with the latter. It was revised in 2005 to become the PPI-R and now comprises 154 items organized into eight subscales. The item scores have been found to group into two overarching and largely separate factors (unlike the PCL-R factors), Fearless-Dominance and Impulsive Antisociality, plus a third factor, Coldheartedness, which is largely dependent on scores on the other two. Factor 1 is associated with social efficacy while Factor 2 is associated with maladaptive tendencies. A person may score at different levels on different factors, but the overall score indicates the extent of psychopathic personality.
Triarchic Psychopathy Measure
The Triarchic Psychopathy Measure, otherwise known as the TriPM, is a 58-item, self-report assessment that measures psychopathy within the three traits identified in the triarchic model: boldness, meanness, and disinhibition. Each trait is measured on separate subscales and added up resulting in a total psychopathy score.
The TriPM includes various components of other measures for assessing psychopathy, including meanness and disinhibition patterns within the psychopathic personality. However, there are differing approaches in the measurement of the boldness construct. The boldness construct is used to highlighting the social and interpersonal implications of the psychopathic personality.
DSM and ICD
There are currently two widely established systems for classifying mental disorders—the International Classification of Diseases (ICD) produced by the World Health Organization (WHO) and the Diagnostic and Statistical Manual of Mental Disorders (DSM) produced by the American Psychiatric Association (APA). Both list categories of disorders thought to be distinct types, and have deliberately converged their codes in recent revisions so that the manuals are often broadly comparable, although significant differences remain.
The first edition of the DSM in 1952 had a section on sociopathic personality disturbances, then a general term that included such things as homosexuality and alcoholism as well as an "antisocial reaction" and "dyssocial reaction". The latter two eventually became antisocial personality disorder (ASPD) in the DSM and dissocial personality disorder in the ICD. Both manuals have stated that their diagnoses have been referred to, or include what is referred to, as psychopathy or sociopathy, although neither diagnostic manual has ever included a disorder officially titled as such.
Other tools
There are some traditional personality tests that contain subscales relating to psychopathy, though they assess relatively non-specific tendencies towards antisocial or criminal behavior. These include the Minnesota Multiphasic Personality Inventory (Psychopathic Deviate scale), California Psychological Inventory (Socialization scale), and Millon Clinical Multiaxial Inventory Antisocial Personality Disorder scale. There is also the Levenson Self-Report Psychopathy Scale (LSRP) and the Hare Self-Report Psychopathy Scale (HSRP), but in terms of self-report tests, the PPI/PPI-R has become more used than either of these in modern psychopathy research on adults.
Comorbidity
Studies suggest strong comorbidity between psychopathy and antisocial personality disorder. Among numerous studies, positive correlations have also been reported between psychopathy and histrionic, narcissistic, borderline, paranoid, and schizoid personality disorders, panic and obsessive–compulsive disorders, but not neurotic disorders in general, schizophrenia, or depression.
Factor 1 and the boldness scale of psychopathy measurements are associated with narcissism and histrionic personality disorder. This is due to a psychopath's cognitive and affective egocentrism. However, while a narcissistic individual might view themselves as confident, they might seek out validation and attention from others to validate their self-worth, whereas a psychopathic individual usually lacks such ambitions.
Attention deficit hyperactivity disorder (ADHD) is known to be highly comorbid with conduct disorder (a theorized precursor to ASPD), and may also co-occur with psychopathic tendencies. This may be explained in part by deficits in executive function. Anxiety disorders often co-occur with ASPD, and contrary to assumptions, psychopathy can sometimes be marked by anxiety; this appears to be related to items from Factor 2 but not Factor 1 of the PCL-R. Psychopathy is also associated with substance use disorders.
Michael Fitzgerald suggested overlaps between (primary) psychopathy and Asperger syndrome in terms of fearlessness, planning of acts, empathy deficits, callous behaviour, and sometimes superficial charisma. Studies investigating similarities and differences between psychopathy and autism indicate that autism and psychopathy are not part of the same construct. Rather both conditions might co-occur in some individuals. Recent studies indicate that some individuals with an autism diagnosis also show callous and unemotional traits (a risk-factor for developing psychopathy), but are less strongly associated with conduct problems. Likewise, some people with Asperger syndrome have shown correlations with the "unemotional" factor and "behavioural dyscontrol" factor of psychopathy, but not the "interpersonal" factor.
It has been suggested that psychopathy may be comorbid with several other conditions than these, but limited work on comorbidity has been carried out. This may be partly due to difficulties in using inpatient groups from certain institutions to assess comorbidity, owing to the likelihood of some bias in sample selection.
Sex differences
Research on psychopathy has largely been done on men and the PCL-R was developed using mainly male criminal samples, raising the question of how well the results apply to women. Men score higher than women on both the PCL-R and the PPI and on both of their main scales. The differences tend to be somewhat larger on the interpersonal-affective scale than on the antisocial scale. Most but not all studies have found broadly similar factor structure for men and women.
Many associations with other personality traits are similar, although in one study the antisocial factor was more strongly related to impulsivity in men and more strongly related to openness to experience in women. It has been suggested that psychopathy in men manifests more as an antisocial pattern while in women it manifests more as a histrionic pattern. Studies on this have shown mixed results. PCL-R scores may be somewhat less predictive of violence and recidivism in women. On the other hand, psychopathy may have a stronger relationship with suicide and possibly internalizing symptoms in women. A suggestion is that psychopathy manifests more as externalizing behaviors in men and more as internalizing behaviors in women. Furthermore, one study has suggested substantial gender differences were found in the etiology of psychopathy. For girls, 75% of the variance in severe callous and unemotional traits was attributable to environmental factors and just 0% of the variance was attributable to genetic factors. In boys, the link was reversed.
Studies have also found that women in prison score significantly lower on psychopathy than men, with one study reporting only 11 percent of violent females in prison met the psychopathy criteria in comparison to 31 percent of violent males. Other studies have also indicated that high psychopathic females are rare in forensic settings.
Management
Clinical
Psychopathy has often been considered untreatable. Its unique characteristics make it among the most refractory of personality disorders, a class of mental illnesses that are already traditionally considered difficult to treat. People with psychopathy are generally unmotivated to seek treatment for their condition, and can be uncooperative in therapy. Attempts to treat psychopathy with the current tools available to psychiatry have been disappointing. Harris and Rice's Handbook of Psychopathy says that there is currently little evidence for a cure or effective treatment for psychopathy; as yet, no pharmacological therapies are known to or have been trialed for alleviating the emotional, interpersonal and moral deficits of psychopathy, and patients with psychopathy who undergo psychotherapy might gain the skills to become more adept at the manipulation and deception of others and be more likely to commit crime. Some studies suggest that punishment and behavior modification techniques are ineffective at modifying the behavior of psychopathic individuals as they are insensitive to punishment or threat. These failures have led to a widely pessimistic view on its treatment prospects, a view that is exacerbated by the little research being done into psychopathy compared to the efforts committed to other mental illnesses, which makes it more difficult to gain the understanding of this condition that is necessary to develop effective therapies.
Although the core character deficits of highly psychopathic individuals are likely to be highly incorrigible to the currently available treatment methods, the antisocial and criminal behavior associated with it may be more amenable to management, the management of which being the main aim of therapy programs in correctional settings. It has been suggested that the treatments that may be most likely to be effective at reducing overt antisocial and criminal behavior are those that focus on self-interest, emphasizing the tangible, material value of prosocial behavior, with interventions that develop skills to obtain what the patient wants out of life in prosocial rather than antisocial ways. To this end, various therapies have been tried with the aim of reducing the criminal activity of incarcerated offenders with psychopathy, with mixed success. As psychopathic individuals are insensitive to sanction, reward-based management, in which small privileges are granted in exchange for good behavior, has been suggested and used to manage their behavior in institutional settings.
Psychiatric medications may also alleviate co-occurring conditions sometimes associated with psychopathy or with symptoms such as aggression or impulsivity, including antipsychotic, antidepressant or mood-stabilizing medications, although none have yet been approved by the FDA for this purpose. For example, a study found that the antipsychotic clozapine may be effective in reducing various behavioral dysfunctions in a sample of high-security hospital inpatients with antisocial personality disorder and psychopathic traits. However, research into the pharmacological treatment of psychopathy and the related condition antisocial personality disorder is minimal, with much of the knowledge in this area being extrapolations based on what is known about pharmacology in other mental disorders.
Legal
The PCL-R, the PCL:SV, and the PCL:YV are highly regarded and widely used in criminal justice settings, particularly in North America. They may be used for risk assessment and for assessing treatment potential and be used as part of the decisions regarding bail, sentence, which prison to use, parole, and whether a youth should be tried as a juvenile or as an adult. There have been several criticisms against its use in legal settings. They include the general criticisms against the PCL-R, the availability of other risk assessment tools that may have advantages, and the excessive pessimism surrounding the prognosis and treatment possibilities of those who are diagnosed with psychopathy.
The interrater reliability of the PCL-R can be high when used carefully in research but tends to be poor in applied settings. In particular Factor 1 items are somewhat subjective. In sexually violent predator cases the PCL-R scores given by prosecution experts were consistently higher than those given by defense experts in one study. The scoring may also be influenced by other differences between raters. In one study it was estimated that of the PCL-R variance, about 45% was due to true offender differences, 20% was due to which side the rater testified for, and 30% was due to other rater differences.
To aid a criminal investigation, certain interrogation approaches may be used to exploit and leverage the personality traits of suspects thought to have psychopathy and make them more likely to divulge information.
United Kingdom
The PCL-R score cut-off for a label of psychopathy is 25 out of 40 in the United Kingdom, instead of 30 as it is in the United States.
In the United Kingdom, "psychopathic disorder" was legally defined in the Mental Health Act (UK), under MHA1983, as "a persistent disorder or disability of mind (whether or not including significant impairment of intelligence) which results in abnormally aggressive or seriously irresponsible conduct on the part of the person concerned". This term was intended to reflect the presence of a personality disorder in terms of conditions for detention under the Mental Health Act 1983. Amendments to MHA1983 within the Mental Health Act 2007 abolished the term "psychopathic disorder", with all conditions for detention (e.g. mental illness, personality disorder, etc.) encompassed by the generic term of "mental disorder".
In England and Wales, the diagnosis of dissocial personality disorder is grounds for detention in secure psychiatric hospitals under the Mental Health Act if they have committed serious crimes, but since such individuals are disruptive to other patients and not responsive to usual treatment methods this alternative to traditional incarceration is often not used.
United States
"Sexual psychopath" laws
Starting in the 1930s, before some modern concepts of psychopathy were developed, "sexual psychopath" laws, the term referring broadly to mental illness, were introduced by some states, and by the mid-1960s more than half of the states had such laws. Sexual offenses were considered to be caused by underlying mental illnesses, and it was thought that sex offenders should be treated, in agreement with the general rehabilitative trends at this time. Courts committed sex offenders to a mental health facility for community protection and treatment.
Starting in 1970, many of these laws were modified or abolished in favor of more traditional responses such as imprisonment due to criticism of the "sexual psychopath" concept as lacking scientific evidence, the treatment being ineffective, and predictions of future offending being dubious. There were also a series of cases where persons treated and released committed new sexual offenses. Starting in the 1990s, several states have passed sexually dangerous person laws, including registration, housing restrictions, public notification, mandatory reporting by health care professionals, and civil commitment, which permits indefinite confinement after a sentence has been completed. Psychopathy measurements may be used in the confinement decision process.
Prognosis
The prognosis for psychopathy in forensic and clinical settings is quite poor, with some studies reporting that treatment may worsen the antisocial aspects of psychopathy as measured by recidivism rates, though it is noted that one of the frequently cited studies finding increased criminal recidivism after treatment, a 2011 retrospective study of a treatment program in the 1960s, had several serious methodological problems and likely would not be approved of today. However, some relatively rigorous quasi-experimental studies using more modern treatment methods have found improvements regarding reducing future violent and other criminal behavior, regardless of PCL-R scores, although none were randomized controlled trials. Various other studies have found improvements in risk factors for crime such as substance abuse. No study has yet examined whether the personality traits that form the core character disturbances of psychopathy could be changed by such treatments.
Frequency
A 2008 study using the PCL:SV found that 1.2% of a US sample scored 13 or more out of 24, indicating "potential psychopathy". The scores correlated significantly with violence, alcohol use, and lower intelligence. A 2009 British study by Coid et al., also using the PCL:SV, reported a community prevalence of 0.6% scoring 13 or more. However, if the scoring was adjusted to the recommended 18 or more, this would have left the prevalence closer to 0.1%. The scores correlated with younger age, male gender, suicide attempts, violence, imprisonment, homelessness, drug dependence, personality disorders (histrionic, borderline and antisocial), and panic and obsessive–compulsive disorders.
Psychopathy has a much higher prevalence in the convicted and incarcerated population, where it is thought that an estimated 15–25% of prisoners qualify for the diagnosis. A study on a sample of inmates in the UK found that 7.7% of the inmates interviewed met the PCL-R cut-off of 30 for a diagnosis of psychopathy. A study on a sample of inmates in Iran using the PCL:SV found a prevalence of 23% scoring 18 or more. A study by Nathan Brooks from Bond University found that around one in five corporate bosses display clinically significant psychopathic traits - a proportion similar to that among prisoners.
Society and culture
In the workplace
There is limited research on psychopathy in the general work populace, in part because the PCL-R includes antisocial behavior as a significant core factor (obtaining a PCL-R score above the threshold is unlikely without having significant scores on the antisocial-lifestyle factor) and does not include positive adjustment characteristics, and most researchers have studied psychopathy in incarcerated criminals, a relatively accessible population of research subjects.
However, psychologists Fritzon and Board, in their study comparing the incidence of personality disorders in business executives against criminals detained in a mental hospital, found that the profiles of some senior business managers contained significant elements of personality disorders, including those referred to as the "emotional components", or interpersonal-affective traits, of psychopathy. Factors such as boldness, disinhibition, and meanness as defined in the triarchic model, in combination with other advantages such as a favorable upbringing and high intelligence, are thought to correlate with stress immunity and stability, and may contribute to this particular expression. Such individuals are sometimes referred to as "successful psychopaths" or "corporate psychopaths" and they may not always have extensive histories of traditional criminal or antisocial behavior characteristic of the traditional conceptualization of psychopathy. Robert Hare claims that the prevalence of psychopathic traits is higher in the business world than in the general population, reporting that while about 1% of the general population meet the clinical criteria for psychopathy, figures of around 3–4% have been cited for more senior positions in business. Hare considers newspaper tycoon Robert Maxwell to have been a strong candidate as a "corporate psychopath".
Academics on this subject believe that although psychopathy is manifested in only a small percentage of workplace staff, it is more common at higher levels of corporate organizations, and its negative effects (for example, increased bullying, conflict, stress, staff turnover, absenteeism, reduction in productivity) often causes a ripple effect throughout an organization, setting the tone for an entire corporate culture. Employees with the disorder are self-serving opportunists, and may disadvantage their own organizations to further their own interests. They may be charming to staff above their level in the workplace hierarchy, aiding their ascent through the organization, but abusive to staff below their level, and can do enormous damage when they are positioned in senior management roles. Psychopathy as measured by the PCL-R is associated with lower performance appraisals among corporate professionals. The psychologist Oliver James identifies psychopathy as one of the dark triadic traits in the workplace, the others being narcissism and Machiavellianism, which, like psychopathy, can have negative consequences.
According to a study from the University of Notre Dame published in the Journal of Business Ethics, psychopaths have a natural advantage in workplaces overrun by abusive supervision, and are more likely to thrive under abusive bosses, being more resistant to stress, including interpersonal abuse, and having less of a need for positive relationships than others.
In fiction
Characters with psychopathy or sociopathy are some of the most notorious characters in film and literature, but their characterizations may only vaguely or partly relate to the concept of psychopathy as it is defined in psychiatry, criminology, and research. The character may be identified as having psychopathy within the fictional work itself, by its creators, or from the opinions of audiences and critics, and may be based on undefined popular stereotypes of psychopathy. Characters with psychopathic traits have appeared in Greek and Roman mythology, Bible stories, and some of Shakespeare's works.
Such characters are often portrayed in an exaggerated fashion and typically in the role of a villain or antihero, where the general characteristics and stereotypes associated with psychopathy are useful to facilitate conflict and danger. Because the definitions, criteria, and popular conceptions throughout its history have varied over the years and continue to change even now, many of the characters characterized as psychopathic in notable works at the time of publication may no longer fit the current definition and conception of psychopathy. There are several archetypal images of psychopathy in both lay and professional accounts which only partly overlap and can involve contradictory traits: the charming con artist, the deranged serial killer and mass murderer, the callous and scheming businessperson, and the chronic low-level offender and juvenile delinquent. The public concept reflects some combination of fear of a mythical bogeyman, the disgust and intrigue surrounding evil, and fascination and sometimes perhaps envy of people who might appear to go through life without attachments and unencumbered by guilt, anguish or insecurity.
| Biology and health sciences | Mental disorder | null |
12408673 | https://en.wikipedia.org/wiki/African%20bullfrog | African bullfrog | The African bullfrog (Pyxicephalus adspersus), also known as the giant bullfrog or the South African burrowing frog, is a species of frog in the family Pyxicephalidae. It is also known as the pixie frog due to its scientific name. It is found in Angola, Botswana, Kenya, Malawi, Mozambique, Namibia, South Africa, Tanzania, Uganda, Zambia, Zimbabwe, and possibly the Democratic Republic of the Congo. It has been extirpated from Eswatini. It has long been confused with the edible bullfrog (P. edulis), and species boundaries between them, including exact range limits, are not fully understood. Additionally, P. angusticeps of coastal East Africa only was revalidated as a separate species in 2013.
The natural habitat of the African bullfrog is moist to dry savanna, subtropical to tropical dry shrubland, intermittent freshwater lakes, intermittent freshwater marshes, arable land, pastureland, canals, and flooded ditches. It is among the largest anurans on the planet, sixth only to the goliath frog, the helmeted water toad, the Lake Junin frog, the Blyth's river frog, and the cane toad.
Appearance
The males weigh up to and grow to in snout–to–vent length. Females are half the size of males, which is unusual, considering that most female amphibians are (generally) larger than males, to help the amplexus. The dorsum is dark olive-green in adults. Juveniles have bright green coloration with a yellow stripe down the back. The abdomen is white to creamy-yellow with breeding males having yellow forelimbs.
Feeding and habits
The African bullfrog is a voracious carnivore, eating insects and other invertebrates, small rodents, reptiles, small birds, fish, and other amphibians that can fit in their mouths. It is also a cannibalistic species—the male African bullfrog is known for occasionally eating the tadpoles he guards, and juveniles also eat tadpoles. An African bullfrog kept at the Pretoria Zoo in South Africa once ate 17 juvenile Rinkhals snakes (Hemachatus haemachatus).
When exposed to dry conditions, they become dormant and may form a cocoon which covers the entire body surface except the external nostrils. The cocoon significantly decreases the rate of evaporative water loss. When it rains, the water softens the cocoon and the frog crawls out.
Reproduction
Breeding typically starts after about of rain over the course of two days. They breed in shallow, temporary water bodies, such as pools, pans, and ditches. Eggs are laid in the shallow edge of the pond, but fertilization takes place above water.
The African bullfrog males call out during the rainy season. The call lasts about a second and can be described as a low-pitched whoop. Males have two breeding strategies, depending on their age. Young males congregate in a small area, perhaps only of shallow water. The larger males occupy the centre of these breeding arenas or leks, and attempt to chase off other males. Often, they fight, causing injury or even killing one another. The dominant male attempts to prevent other males from breeding. A female approaches the group of males by swimming along at the surface until she is within a few metres of the group. Then, the female dives to avoid the smaller males and surfaces in the defended area of a larger male in the middle of the group. This helps to ensure that she mates with the dominant male.
The female lays about 3,000 to 4,000 eggs at a time. The tadpoles hatch, and after two days, start feeding on vegetation, small fish, invertebrates, and even each other. Defending males continue to watch over the tadpoles, which metamorphose within three weeks. During the tadpole's development, the father guards his young. Due to the male bullfrog's overprotective behaviour, he pounces and bites anything that he views as a threat. If the pool is in danger of drying out, the father uses his legs and head to dig a canal from the drying pond to a bigger pond. He continues to guard the tadpoles until they are old enough to fend for themselves, although he may also eat some of them.
Status
The African bullfrog is listed as "least concern" globally. Before 2013, it was considered "near threatened" in South Africa due to habitat loss. However, conservation initiatives, such as habitat restoration, awareness campaigns, and legal protections, have been instrumental in reversing this trend, resulting in its reclassification to "least concern" in July 2013.
Human use
The African bullfrog is an exotic pet in many countries around the world. Animals sold are generally bred in captivity. Pet African bullfrogs may live for 35 years in captivity. As pets, African bullfrogs are considered to be easygoing and low-maintenance in terms of their care.
It is considered a delicacy in Namibia.
Gallery
| Biology and health sciences | Frogs and toads | Animals |
3293252 | https://en.wikipedia.org/wiki/Psychodidae | Psychodidae | Psychodidae, also called drain flies, sink flies, filter flies, sewer flies, or sewer gnats, is a family of true flies. Some genera have short, hairy bodies and wings, giving them a "furry" moth-like appearance, hence one of their common names, moth flies. Members of the sub-family Phlebotominae, which are hematophagous (feed on blood), may be called sand flies in some countries, although this term is also used for other unrelated flies, such as horse flies (Tabanidae) and no-see-ums (Ceratopogonidae).
There are more than 2,600 described species worldwide, most of them native to the humid tropics. This makes them one of the most diverse families of their order. Drain flies sometimes inhabit plumbing drains and sewage systems, where they are harmless, but may be a persistent annoyance.
Life cycle
The larvae of the subfamilies Psychodinae, Sycoracinae and Horaiellinae live in aquatic to semi-terrestrial or sludge-based habitats, including bathroom sinks, where they feed on bacteria and can become problematic. The larvae of the most commonly encountered species are nearly transparent with a non-retractable black head and can sometimes be seen moving along the moist edges of crevices in shower stalls or bathtubs or submerged in toilet water. The larval form of the fly is usually between long, and is shaped like a long, thin, somewhat flattened cylinder. The body lacks prolegs, but the body segments are divided into a series of rings called annuli (singular is annulus). Some of these rings will have characteristic plates on the dorsal side. The larval thorax is not significantly larger than the abdomen, giving the larvae a more "worm-like" appearance than those of most aquatic insects.
In some species, the larvae can secure themselves to surfaces of their environment using "attachment disks" on their ventral side. Like mosquito larvae, they cannot absorb oxygen through water, and instead breathe via a small dark tube (a spiracle) on their posterior end — they must regularly reach the surface to obtain oxygen. The larval stage lasts for between 9 and 15 days, depending on species, temperature, and environment. There are four instar stages.
In small numbers, the larvae are sometimes considered beneficial, as their strong jaws can cut through the hair and sludge waste in drains which might otherwise form clogs. However, unless this sludge layer is removed entirely, the adult flies will continue to find it and lay more eggs.
While the biting midges also have larvae that have no prolegs and which also have attachment disks, the larvae of the netwinged midges can be distinguished from those of the moth fly by the multiple deep lateral constrictions of the latter.
The pupal stage lasts between 20 and 40 hours. During this stage, the insect does not feed, but stays submerged near the water surface, still breathing through a spiracle, and soon metamorphoses into an adult fly, which bursts through a seam in the pupal casing and emerges onto the water's surface.
The adults are half as long as the larvae, but are much broader in appearance, with a pair of hairy wings held pitched-roof-like over the body. The wings have the most elementary venation of any of the Diptera, having little more than a series of parallel veins without crossveins.
The adults are typically nocturnal, though they orient themselves around lights and may appear to be attracted to light and odors. They are erratic fliers, and are often seen walking or running rapidly as well as taking flight. They are most active at night, but may also be seen during daylight, or near windows, lights, or illuminated display panels.
The adults live for about 20 days, during which they will breed only once, often within hours of emerging from their pupal casings. Females will lay their eggs (between 30 and 100) just above the water line inside moist drains. Within 48 hours these eggs hatch into drain worms, the larval form.
Health effects
The drain flies which are commonly found in bathrooms, Clogmia albipunctata, are not known to carry any human diseases, but have been known to be an opportunistic agent of myiasis. However, the subfamily of Phlebotominae does feed on blood with the ability to transmit (tropical) diseases, and Sycorax silacea can transmit microfilaria. Inhalation of insect fragments may cause asthma.
Taxonomy
This family has seven subfamilies that contain more than 2600 described species.
Horaiellinae Enderlein, 1937
Horaiella Tonnoir, 1933
Protohoraiella Curler, Krzeminski & Skibinska, 2019 Burmese amber, Myanmar Late Cretaceous (Cenomanian)
Bruchomyiinae Alexander, 1921 - selected genera:
Alexanderia Wagner & Kvifte, 2018 (Oriental)
Boreofairchildia Wagner & Stuckenberg, 2016 (Americas)
Bruchomyia Alexander, 1921 (South America):
Eutonnoiria Alexander, 1940 (Central Africa)
Laurenceomyia Wagner & Stuckenberg, 2016 (South America)
Nemopalpus Macquart, 1838
Notofairchildia Wagner & Stuckenberg, 2016
Hoffeinsodes Wagner, 2017 Baltic amber, Eocene
Palaeoglaesum Wagner, 2017 Burmese amber, Myanmar
Phlebotominae Rondani, 1840
Australophlebotomus Theodor, 1948
Bichromomyia Artemiev, 1991
Brumptomyia França & Parrot, 1921 (Mexico to South America)
Chinius Leng, 1985 (2 species: China, Thailand)
Dampfomyia Addis, 1945
Datzia Stebner et al., 2015 (Burmese amber, Cenomanian)
Deanemyia Galati, 1995
Evandromyia Mangabeira, 1941
Edentomyia Galati, Andrade-Filho, da Silva & Falcão, 2003 (Brazil)
Expapillata Galati, 1995
Hertigia Fairchild, 1949
Idiophlebotomus Quate & Fairchild, 1961
Libanophlebotomus Azar et al., 1999 Lebanese amber Early Cretaceous (Barremian)
Lutzomyia França, 1924 (North and South America)
Mandalayia Stebner et al., 2015 (Burmese amber, Cenomanian)
Martinsmyia Galati, 1995
Mesophlebotomites Azar et al., 1999 Lebanese amber, Barremian
Micropygomyia Barretto, 1962
Migonemyia Galati, 1995
Nyssomyia Barretto, 1962
Oligodontomyia Galati, 1995
Palaeomyia Poinar 2004 Burmese amber, Albian
Phlebotomites Stebner et al., 2015 Lebanese amber, Barremian, Burmese amber, Cenomanian
Phlebotoiella Solórzano Kraemer and Wagner 2009 Cambay amber, India, Eocene
Phlebotomus Rondani& Berté, 1840 (Europe, Africa, Asia, Australia)
Pintomyia Costa Lima, 1932
Pressatia Mangabeira, 1942
Protopsychodinae Stebner et al., 2015
Protopsychoda Azar et al., 1999 Lebanese amber, Barremian
Psathyromyia Barretto, 1962
Psychodopygus Mangabeira, 1941
Sciopemyia Barretto, 1962
Sergentomyia França & Parrot, 1920 (Europe, Africa, Asia, Australia)
Trichophoromyia Barretto, 1962
Viannamyia Mangabeira, 1941
Warileya Hertig, 1948 (Central and South America)
Protopsychodinae Stebner et al., 2015
Datzia Stebner et al., 2015 (Burmese amber, Cenomanian)
Mandalayia Stebner et al., 2015 (Burmese amber, Cenomanian)
Protopsychoda Azar et al., 1999 Lebanese amber, Barremian
Psychodinae Newman, 1834
Abcharis Tkoc and Jezek, 2013 (= Notiocharis Eaton, 1913, preoccupied) (Australia)
Alloeodidicrum Duckhouse, 1990 (Australia)
Arisemus Satchell, 1955
Atrichobrunettia Satchell, 1953
Australopericoma Vaillant, 1975
Balbagathis Quate, 1996
Bazarella Vaillant, 1961
Berdeniella Vaillant, 1976
Boreoclytocerus Duckhouse, 1978
Breviscapus Quate, 1955
Brunettia Annandale, 1910
Clogmia Enderlein, 1937
Clytocerus Eaton, 1904
Didicrum Enderlein, 1937
Epacretron Quate, 1965
Eremolobulosa Duckhouse, 1990 (Australia)
Eurygarka Quate, 1959
Feuerborniella Vaillant, 1974
Gerobrunettia Quate & Quate, 1967
Lepimormia Enderlein, 1937
Lepidiella Enderlein, 1937
Lobulosa Szabo, 1960
Maruina Müller, 1895 (Americas)
Matuna Stebner and Solórzano Kraemer 2014 Mexican amber, Miocene
Megapsychoda Azar and Nel 2002 Crato Formation, Brazil, Early Cretaceous (Aptian)
Mormia Enderlein, 1937
Neoarisemus Botosaneanu & Vaillant, 1970
Paleopsychoda Azar et al., 1999 Lebanese amber, Barremian, Jordanian amber, Albian, Taimyr amber, Russia, Albian
Panimerus Eaton, 1913
Paralibanopsychoda Azar and Nel 2002 Lebanese amber, Barremian
Paramormia Enderlein, 1937
Parasetomima Duckhouse, 1968 (South America)
Paratelmatoscopus Satchell, 1953 (Australia)
Pericoma Haliday, in Walker, 1856
Peripsychoda Enderlein, 1937
Philosepedon Eaton, 1904 (Europe, North and Central America)
Pneumia Enderlein, 1937 (= Satchelliella Vaillant, 1979)
Psychoda Latreille, 1796
Rotundopteryx Duckhouse, 1990 (Australia)
Saraiella Vaillant, 1981
Setomima Enderlein, 1937
Stupkaiella Vaillant, 1973
Succinarisemus Wagner, 2002 Mexican amber, Dominican amber, Miocene
Szaboiella Vaillant, 1979
Telmatoscopus Eaton, 1904
Thornburghiella Vaillant, 1982
Threticus Eaton, 1904
Tinearia Schellenberg, 1803
Tonnoiriella Vaillant, 1982
Trichopsychoda Tonnoir, 1922
Ulomyia Walker, 1856 (= Saccopterix Haliday, in Curtis, 1839, preoccupied)
Vaillantodes Wagner, 2002 ( = Vaillantia Wagner, 1993, preoccupied)
Wightipsychoda Azar 2019 Bembridge Marls, United Kingdom, Priabonian
Sycoracinae Jung, 1954
Aposycorax Duckhouse, 1972
Palaeoparasycorax Stebner et al., 2015 (Burmese amber, Cenomanian)
Parasycorax Duckhouse, 1972
Sycorax Haliday, in Curtis, 1839
Trichomyiinae Tonnoir, 1922
Axenotrichomyia Azar et al., 2015 Burmese amber, Cenomanian
Eatonisca Meunier, 1905 Baltic, Bitterfeld amber, Eocene
Eotrichomyia Meunier Oise amber, France Eocene (Ypresian)
Trichomyia Haliday, in Curtis, 1839
Xenotrichomyia Azar et al., 2015 New Jersey amber, Late Cretaceous (Turonian)
Incertae sedis
Bamara Stebner et al., 2015 Burmese amber, Cenomanian
Cretapsychoda Azar et al., 1999 Lebanese amber, Barremian
Eochaoborites Hong, 2002 Fushun amber, China, Ypresian
Eophlebotomus Cockerell, 1920 Lebanese amber, Barremian, Charentese amber, France, Cenomanian, Burmese amber, Cenomanian
Liassopsychodina Ansorge, 1994 Green Series, Germany, Early Jurassic (Toarcian)
Libanopsychoda Azar et al., 1999 Lebanese amber, Barremian
Mesopsychoda Brauer et al., 1889 Cheremkhovskaya Formation, Russia, Toarcian
Protopsychoda Azar et al., 1999 Lebanese amber, Barremian
Tanypsycha Ansorge, 1994 Green Series, Germany, Toarcian
Triassopsychoda Blagoderov and Grimaldi. 2007 Cow Branch Formation, North Carolina, Late Triassic (Norian)
Xenopsychoda Azar and Ziadé, 2005 Lebanese amber, Barremian
| Biology and health sciences | Flies (Diptera) | Animals |
3298595 | https://en.wikipedia.org/wiki/Carcinus%20maenas | Carcinus maenas | Carcinus maenas is a common littoral crab. It is known by different names around the world. In the British Isles, it is generally referred to as the shore crab, or green shore crab. In North America and South Africa, it bears the name European green crab.
C. maenas is a widespread invasive species, listed among the 100 of the World's Worst Invasive Alien Species. It is native to the north-east Atlantic Ocean and Baltic Sea, but has colonised similar habitats in Australia, South Africa, South America and both Atlantic and Pacific coasts of North America. It grows to a carapace width of , and feeds on a variety of mollusks, worms, and small crustaceans, affecting a number of fisheries. Its successful dispersal has occurred by a variety of mechanisms, such as on ships' hulls, sea planes, packing materials, and bivalves moved for aquaculture.
Description
C. maenas has a carapace up to long and wide, but can be larger outside its native range, reaching wide in British Columbia. The carapace has five short teeth along the rim behind each eye, and three undulations between the eyes. The undulations, which protrude beyond the eyes, are the simplest means of distinguishing C. maenas from the closely related C. aestuarii, which can also be an invasive species. In C. aestuarii, the carapace lacks any bumps and extends forward beyond the eyes. Another characteristic for distinguishing the two species is the form of the first and second pleopods (collectively the gonopods), which are straight and parallel in C. aestuarii, but curve outwards in C. maenas.
The colour of C. maenas varies greatly, from green to brown, grey, or red. This variation has a genetic component, but is largely due to local environmental factors. In particular, individuals which delay moulting become red-coloured rather than green. Red individuals are stronger and more aggressive, but are less tolerant of environmental stresses, such as low salinity or hypoxia. Juvenile crabs on average display greater patterning than adults.
Native and introduced range
C. maenas is native to European and North African coasts as far as the Baltic Sea in the east, and Iceland and Central Norway in the north, and is one of the most common crabs throughout much of its range. In the Mediterranean Sea, it is replaced by the closely related Mediterranean Green Crab species Carcinus aestuarii.
C. maenas was first observed on the east coast of North America in Massachusetts in 1817, and may now be found from South Carolina northwards; by 2007, this species had extended its range northwards to Placentia Bay, Newfoundland. In 1989, the species was found in San Francisco Bay, California, on the Pacific Coast of the United States. Until 1993, it was not able to extend its range, but reached Oregon in 1997, Washington in 1998, and British Columbia in 1999, thus extending its range by in 10 years. they were just south of Alaska, and were expected to enter Alaska next. By 2003, C. maenas had extended to South America with specimens discovered in Patagonia.
In Australia, C. maenas was first reported "in the late 1800s" in Port Phillip Bay, Victoria, although the species was probably introduced as early as the 1850s. It has since spread along the south-eastern and south-western seaboards, reaching New South Wales in 1971, South Australia in 1976 and Tasmania in 1993. One specimen was found in Western Australia in 1965, but no further discoveries have been reported in the area since.
C. maenas first reached South Africa in 1983, in the Table Docks area near Cape Town. Since then, it has spread at least as far as Saldanha Bay in the north and Camps Bay in the south, over apart.
Appearances of C. maenas have been recorded in Brazil, Panama, Hawaii, Madagascar, the Red Sea, Pakistan, Sri Lanka, and Myanmar; however, these have not resulted in invasions, but remain isolated findings. Japan has been invaded by a related crab, either Carcinus aestuarii or a hybrid of Carcinus aestuarii and C. maenas.
Based on the ecological conditions, C. maenas could eventually extend its range to colonise the Pacific Coast of North America from Baja California to Alaska. Similar ecological conditions are to be found on many of the world's coasts, with the only large potential area not to have been invaded yet being New Zealand; the New Zealand government has taken action, including the release of a Marine Pest Guide in an effort to prevent colonisation by C. maenas.
In 2019 C. maenas was first found in Lummi Bay, Lummi Indian Reservation, Whatcom County, Washington, USA. The Nation began trapping and removing the crabs in an effort to get rid of them. Then in 2020 hundreds were found in traps, and it became clear that more intensive trapping will be necessary to keep their numbers down. Eradication will not be possible.
Over a 19-year study concluding in 2020, Oregon's Coos Bay was found to have an established and increasing population.
While in 2020 less than 3,000 were trapped, trapping yielded >79,000 in 2021. This led the Lummi Indian Business Council to declare a disaster in November 2021 and the Washington Department of Fish & Wildlife to request emergency funding from the Governor.
Ecology
C. maenas can live in all types of protected and semiprotected marine and estuarine habitats, including those with mud, sand, or rock substrates, submerged aquatic vegetation, and emergent marsh, although soft bottoms are preferred. C. maenas is euryhaline, meaning that it can tolerate a wide range of salinities (from 4 to 52 ‰), and survive in temperatures of . The wide salinity range allows C. maenas to survive in the lower salinities found in estuaries, and the wide temperature range allows it to survive in extremely cold climates beneath the ice in winter. A molecular biological study using the COI gene found genetic differentiation between the North Sea and the Bay of Biscay, and even more strongly between the populations in Iceland and the Faroe Islands and those elsewhere. This suggests that C. maenas is unable to cross deeper water.
Females can produce up to 185,000 eggs, and larvae develop offshore in several stages before their final moult to juvenile crabs in the intertidal zone. Young crabs live among seaweeds and seagrasses, such as Posidonia oceanica, until they reach adulthood.
C. maenas has the ability to disperse by a variety of mechanisms, including ballast water, ships' hulls, packing materials (seaweeds) used to ship live marine organisms, bivalves moved for aquaculture, rafting, migration of crab larvae on ocean currents, and the movement of submerged aquatic vegetation for coastal zone management initiatives. C. maenas dispersed in Australia mainly by rare long-distance events, possibly caused by human actions.
C. maenas is a predator, feeding on many organisms, particularly bivalve molluscs (such as clams – up to 40 clams per day, oysters, and mussels), polychaetes, and small crustaceans – including other crabs up to their own size. They are primarily diurnal, although activity also depends on the tide, and crabs can be active at any time of day. In California, preferential predation of C. maenas on native clams (Nutricola spp.) resulted in the decline of the native clams and an increase of a previously introduced clam (the amethyst gem clam, Gemma gemma), although C. maenas also voraciously preys on introduced clams such as Potamocorbula amurensis. The soft-shell clam (Mya arenaria) is a preferred prey species of C. maenas. Consequently, it has been implicated in the destruction of the soft-shell clam fisheries on the east coast of the United States and Canada, and the reduction of populations of other commercially important bivalves (such as scallops, Argopecten irradians, and northern quahogs, Mercenaria mercenaria). The prey of C. maenas includes the young of bivalves and fish, although the effect of its predation on winter flounder, Pseudopleuronectes americanus is minimal. C. maenas can, however, have substantial negative impacts on local commercial and recreational fisheries, by preying on the young of species, such as oysters (adults' shells are too tough for C. maenas to crack) and the Dungeness crab, or competing with them for resources and eating the Zostera marina that Dungeness and juvenile salmon depend upon for habitat. Colder water temperatures reduce overall feeding rates of C. maenas.
To protect itself against predators, C. maenas uses different camouflage strategies depending on their habitat: crabs in mudflats try to resemble their surroundings with colours similar to the mud while crabs in rock pool use disruptive coloration.
Control
Due to its potentially harmful effects on ecosystems, various efforts have been made to control introduced populations of C. maenas around the world. In Edgartown, Massachusetts, a bounty was levied in 1995 for catching C. maenas, to protect local shellfish, and 10 tons were caught.
Some evidence shows that the native blue crab in eastern North America, Callinectes sapidus, is able to control populations of C. maenas; numbers of the two species are negatively correlated, and C. maenas is not found in the Chesapeake Bay, where C. sapidus is most frequent. On the west coast of North America, C. maenas appears to be limited to upper estuarine habitats, in part because of predation by native rock crabs (Romaleon antennarium and Cancer productus) and competition for shelter with a native shore crab, Hemigrapsus oregonensis. Host specificity testing has recently been conducted on Sacculina carcini, a parasitic barnacle, as a potential biological control agent of C. maenas. In the laboratory, Sacculina settled on, infected, and killed native California crabs, including the Dungeness crab, Metacarcinus magister (formerly Cancer magister), and the shore crabs Hemigrapsus nudus, Hemigrapsus oregonensis and Pachygrapsus crassipes. Dungeness crabs were the most vulnerable of the tested native species to settlement and infection by the parasite. Although Sacculina did not mature in any of the native crabs, developing reproductive sacs were observed inside a few M. magister and H. oregonensis crabs. Any potential benefits of using Sacculina to control C. maenas on the west coast of North America would need to be weighed against these potential nontarget impacts.
Use as a food
In its native range, European green crab is mostly used as an ingredient in soups and sauces. However, the closely related Mediterranean green crab (C. aestuarii) has a thriving culinary market in Italy where fishermen known as moecante cultivate soft-shell green crabs ( in Venetian, in Italian) and sell hard-shell crabs for their roe (). Several groups in New England have successfully adapted these methods to produce soft-shell green crabs from the invasive species.
In New England where invasive green crab populations are high, various groups have looked into utilizing green crabs in cuisine. In 2019, The Green Crab Cookbook was released and included recipes for soft-shell green crab, green crab roe, green crab stock, and green crab meat. One of the book's co-authors went on to found Greencrab.org, an organization dedicated to developing culinary markets for the invasive green crab. In addition to partnering with local chefs and wholesalers for supply chain development and market studies, Greencrab.org has continued to develop green crab recipes and processing techniques.
Researchers at the University of Maine have actively been developing value-added green crab products, with the goals of driving business interest, stimulating a commercial green crab fishery, and alleviating predation effects. Specifically, one study evaluated the consumer acceptability of empanadas (fried, stuffed pastries) which contained varying amounts of green crab mince meat. The empanadas were rated between "like slightly" and "like moderately" for overall acceptability by a consumer panel (n=87). Furthermore, about two-thirds of the panelists would "probably" or "definitely" buy the empanadas if available locally. Additionally, the same researchers developed a patty product made from green crab mince meat using restructuring additives (transglutaminase, dried egg white, isolated soy protein). Although a successful green crab patty was developed, the restructuring additives may have had greater functionality in a raw crab meat system, as opposed to the fully cooked mince that was used in the present study. The results from both studies are considered promising, especially considering that these were initial rounds of green crab product development.
In the past, Legal Sea Foods, an East Coast restaurant chain, experimented with green crabs, creating a green crab stock in their test kitchen during the winter of 2015. In June 2022 Tamworth Distilling, a New Hampshire distillery, teamed up with the University of New Hampshire's NH Green Crab Project to develop House of Tamworth Crab Trapper, which is billed as being "made with a bourbon base steeped with a custom crab, corn and spice blend mixture".
Fishery
C. maenas is fished on a small scale in the northeast Atlantic Ocean, with about 1200 tonnes being caught annually, mostly in France and the United Kingdom. In the northwest Atlantic, C. maenas was the subject of fishery in the 1960s, and again since 1996, with up to 86 tonnes being caught annually.
Taxonomic history
Carcinus maenas was first given a binomial name, Cancer maenas, by Carl Linnaeus in his 1758 10th edition of Systema Naturae. An earlier description was published by Georg Eberhard Rumphius in his 1705 work De Amboinsche Rariteitkamer, calling the species Cancer marinus sulcatus, but this antedates the starting point for zoological nomenclature. A number of later synonyms have also been published:
Monoculus taurus Slabber, 1778
Cancer granarius Herbst, 1783
Cancer viridis Herbst, 1783
Cancer pygmaeus Fabricius, 1787
Cancer rhomboidalis Montagu, 1804
Cancer granulatus Nicholls, 1943
Megalopa montagui Leach, 1817
Portunus menoides Rafinesque-Schmaltz, 1817
Portunus carcinoides Kinahan, 1857
The lectotype chosen for the species came from Marstrand, Sweden, but it is assumed to have been lost. In 1814, writing for The Edinburgh Encyclopaedia, William Elford Leach erected a new genus, Carcinus to hold this species alone (making it the type species of the genus, by monotypy). In 1847, Nardo described a distinct subspecies occurring in the Mediterranean Sea, which is now recognised as a distinct species, Carcinus aestuarii.
Neurochemistry
Particular amino acids in particular signaling peptides of C. maenas are protonated by pH changes currently () occurring or likely to be reached in the course of future climate change. This significantly alters peptide structure and peptide-mediated behaviours (brood care and egg ventilation requiring ~10x the normal peptide concentration). The requirement of higher concentration may be due to lowered binding affinity in the sensory epithelium. This effect is very reversible.
Physiochemistry
The usual decrease in extracellular chloride due to increased extracellular bicarbonate is avoided if C. maenas is first acclimated to the increased CO. While this may be due to the already-high extracellular chloride levels in this species, it may instead be because moderately higher CO increases these levels through some unrelated mechanism.
Changes in pH due to sodium and magnesium can alter extracellular iron concentrations.
| Biology and health sciences | Crabs and hermit crabs | Animals |
3300228 | https://en.wikipedia.org/wiki/Herpetic%20gingivostomatitis | Herpetic gingivostomatitis | Gingivostomatitis is a combination of gingivitis and stomatitis, or an inflammation of the oral mucosa and gingiva. Herpetic gingivostomatitis is often the initial presentation during the first ("primary") herpes simplex infection. It is of greater severity than herpes labialis (cold sores) which is often the subsequent presentations. Primary herpetic gingivostomatitis is the most common viral infection of the mouth.
Primary herpetic gingivostomatitis (PHGS) represents the clinically apparent pattern of primary herpes simplex virus (HSV) infection, since the vast majority of other primary infections are symptomless. PHGS is caused predominantly by HSV-1 and affects mainly children. Prodromal symptoms, such as fever, anorexia, irritability, malaise and headache, may occur in advance of disease. The disease presents as numerous pin-head vesicles, which rupture rapidly to form painful irregular ulcerations covered by yellow–grey membranes. Sub-mandibular lymphadenitis, halitosis and refusal to drink are usual concomitant findings.
Signs and symptoms
The symptoms can be mild or severe and may include:
Not able to chew or swallow
Sores on the inside of the cheeks or gums
Fever
General discomfort, uneasiness, or ill feeling
Very sore mouth with no desire to eat
Halitosis (bad breath)
Causes
Herpetic gingivostomatitis is an infection caused by the herpes simplex virus (HSV). The HSV is a double-stranded DNA virus categorised into two types; HSV-1 and HSV-2. HSV-1 is predominantly responsible for oral, facial and ocular infections whereas HSV-2 is responsible for most genital and cutaneous lower herpetic lesions. Both HSV-1, and HSV-2 can be the cause of herpetic gingivostomatitis, although HSV-1 is the source of infection in around 90% of cases.
Herpetic gingivostomatitis infections can present as acute or recurrent. Acute infection refers to the first invasion of the virus, and recurrent is when reactivation of the latent virus occurs. Acute herpetic gingivostomatitis primarily occurs in children, particularly of those under the age of six years old.
On external surfaces the virus is short lived, however it is extremely contagious. Most people acquire the virus via direct contact, it can enter the body by disrupting the integrity of skin, mucous membranes or enter via infected secretions such as saliva. The virus replicates once it has penetrated the epithelial cell, then it travels to the corresponding nerve ganglion (i.e. trigeminal ganglion) via sensory nerves endings. At the nerve ganglion the virus enters a latent phase and remains dormant until it is reactivated. Reactivation can be spontaneous or stimulated by a number of factors such as: reinfection by direct effect of stimuli, immunosuppression, ultraviolet light, febrile illnesses and stress.
Risk factors
Age: Primary herpetic gingivostomatitis is common in children from 6 months to 5 years old. This virus is also common in young adults aged around 20–25.
Immune system: The prevalence and severity of the disease is dependent on the host's immune response and the virulence of the virus.
Environment: As this virus is very contagious it has the potential to spread quickly in enclosed environments e.g. nurseries and orphanages.
Epidemiology: Those living in developing countries are at a higher risk of HSV-1 infection. It has been reported that around a 1/3rd of children living in developing countries are HSV-1 positive by 5 years old and 70-80% of the population are infected by the age of adolescence. In developed countries only 20% of children are infected at the age of 5 and there is no significant increase in disease prevalence until 20–40 years old where the percentage of infected individuals ranges from 40-60%
Socio-economic status: Those with a lower income have a higher risk of HSV-1 infection at a younger age.
Race: Studies have demonstrated that in the USA 35% of African Americans by the age of 5 have presented with the disease whereas only 18% of White Americans are affected.
Pathophysiology
Herpetic gingivostomatitis originates from a primary infection of HSV-1. The series of events that take place during this infection include replication of the herpes simplex virus, cell lysis and finally, destruction of the mucosal tissue.
HSV-1 can very easily enter and replicate within epidermal and dermal cells through skin or mucosal surfaces which have abrasions. This results in numerous small vesicles or blisters of up to 1-2mm on the oral mucosa, erosions on the lips, eventual hemorrhagic crusting and even ulceration, covered by a yellowish-grey pseudomembrane, surrounded by an erythematous halo.
As the virus continues to replicate and incolulate in great amounts, it can enter autonomic or sensory ganglia, where it travels within axons to reach ganglionic nerve bodies. HSV-1 most commonly infects the trigeminal ganglia, where it remains latent. If reactivated, it presents as herpes labialis, also known as cold sores.
Diagnosis
Histopathology
The histological appearance of a herpetic infection on the mucosa includes degeneration of stratified squamous epithelial cells, the loss of intercellular connections and inflammatory infiltrate around the capillaries of the dermis layer. An intact herpetic vesicle presents as an intraepithelial blister histologically. This vesicle is caused by rupture and distension of the virally epithelial cells by intracellular oedema and coalescence of disrupted cells.
Rupturing of the infected cells cause a great number of viral particles to be released, rendering them the ability to affect adjacent epithelial cells and even the sensory axons of the trigeminal nerve. Histologically, these infected cells have an eosinophilic cytoplasm and large, pale vesicular nuclei, appearing swollen under the microscope. The cytoplasms of the infected cells fuse, collectively forming giant cells with many nuclei. The balloon cells and multi-nucleated giant cells can often be identified in smears taken from an intact vesicle or from one which has been recently ruptured.
The lamina propria shows a variable inflammatory infiltrate, the density of which depends on the stage and severity of the disease, and inflammatory cells also extend into the epithelium.
Cowdry type A bodies are intranuclear inclusion bodies visible under light microscopy. They show electron dense glycoproteins and viral capsids. Both Cowdry type A bodies can both be found in varicella zoster and herpetic gingivostomatitis, making it impossible to distinguish between both eosinophilic bodies. One way to distinguish between the herpes virus (and hence herpetic gingivostomatitis) and varicella virus is by direct immunohistochemistry using fluorescent antibodies.
Differential diagnosis
Diagnosis of HG is very much clinically based. Therefore, it is imperative to rule out other diseases that present similarly, bearing in mind the past medical history of the patient.
Some differential diagnoses to bear in mind when considering herpetic gingivostomatitis are:
Teething in infants: A study mentioned that "primary tooth eruption begins at about the time that infants are losing maternal antibody protection against the herpes virus. Also, reports on teething difficulties have recorded symptoms which are remarkably consistent with primary oral herpetic infection such as fever, irritability, sleeplessness, and difficulty with eating." Another study highlighted that "younger infants with higher residual levels of antibodies would experience milder infections and these would be more likely to go unrecognized or be dismissed as teething difficulty."
Herpangina: It is a disease that is caused by the Coxackie A virus rather than a herpes virus. In herpangina, ulcers are usually isolated to the soft palate and anterior pillar of the mouth. In herpetic gingivostomatitis, lesions can be found in these locations, but they are almost always accompanied by ulcerations on the gums, lips, tongue or buccal mucosa and/or by hyperemia, hypertrophy or hemorrhage of the gums.
Hand Foot and Mouth Disease: Similar to herpangina, hand foot and mouth disease occurs predominantly in children. It is caused by Coxsackie A and B virus, and lesions or blisters are found bilaterally on the hands, feet and mouth of the patient.
Oral candidiasis: Also known as thrush, herpetic gingivostomatitis can often be differentiated from these microorganism/bacterial causing white plaques on the palate, buccal mucosa, tongue, oropharynx etc.
Apthous stomatitis: They are commonly known as apthous ulcers, and are characterized by grey membranes and peripheral erythema. Lesions/ulcers for herpetic gingivostomatitis may also be found on the palate and keratinzied gingivae hence aphthous ulcers can be ruled out.
Stevens–Johnson syndrome: Stevens–Johnson syndrome is characterized by early symptoms of malaise and fever, and shortly after that erythema, purpura and plaques on the skin, which often progresses to epidermal necrosis and sloughing in extreme cases.
Infectious mononucleosis - Infectious Mononucleosis presents with a high fever and lymphadenopathy, which is may or may not be presented in the symptoms of herpetic gingivostomatitis. However, upon closer oral examination, ulceration, petechiae and occasional gingivostomatitis may be spotted.
Behcet's syndrome - It is an inflammatory disorder in which the presenting symptoms are recurrent aphthous ulcers, and severe cases may result in the patient having genital lesions, gastro-intestinal problems, and even arthritis.
Varicella - Small ulcers on the back of the oral cavity and vesicular lesions on the scalp and trunk are common for Varicella. It is ruled out as the location of the infections are unilateral, unlike herpetic gingivostomatitis which is bilateral.
Treatment
The aim of treatment is mostly supportive such as pain control, duration of symptoms, viral shedding and in some cases, preventing outbreak. Antibiotics are rarely prescribed to treat bacterial superinfection of oral lesions. Antiviral drugs are used to treat herpetic gingivostomatitis such as aciclovir, valaciclovir, famciclovir, and in resistance cases foscarnet can be used. Treatment does not prevent recurrence. Most individuals who are immunocompetent will fully recover from recurrent herpes labialis in 7 to 14 days. However treatment with antipyretics, oral anaesthetics and analgesics is often needed. In severe cases of herpetic gingivostomatitis, mouth rinses are useful in relieving oral discomfort. These contain topical anaesthetic agents such as lidocaine and diphenhydramine as well as coating agents such as magnesium-containing antacids. In order to prevent dehydration, oral fluid intake is encouraged. Other treatment options include good oral hygiene and gentle debridement of the mouth.
A number of precautions can be taken to reduce the risk of infection of HSV including;
Avoid saliva exchange with those with active HSV
Avoid direct contact with active lesions (if contact occurred, ensure area is washed adequately)
Those with recurrent HSV can apply sunscreen lip balm as ultraviolet light is a stimulus for the infection
Epidemiology
| Biology and health sciences | Viral diseases | Health |
5897742 | https://en.wikipedia.org/wiki/Social%20media | Social media | Social media are interactive technologies that facilitate the creation, sharing and aggregation of content (such as ideas, interests, and other forms of expression) amongst virtual communities and networks. Common features include:
Online platforms that enable users to create and share content and participate in social networking.
User-generated content—such as text posts or comments, digital photos or videos, and data generated through online interactions.
Service-specific profiles that are designed and maintained by the social media organization.
Social media helps the development of online social networks by connecting a user's profile with those of other individuals or groups.
The term social in regard to media suggests platforms enable communal activity. Social media can enhance and extend human networks. Users access social media through web-based apps or custom apps on mobile devices. These interactive platforms allow individuals, communities, and organizations to share, co-create, discuss, participate in, and modify user-generated or self-curated content. Social media is used to document memories, learn, and form friendships. They may be used to promote people, companies, products, and ideas. Social media can be used to consume, publish, or share news.
Popular social media platforms with over 100 million registered users include Twitter, Facebook, WeChat, ShareChat, Instagram, Pinterest, QZone, Weibo, VK, Tumblr, Baidu Tieba, Threads and LinkedIn. Depending on interpretation, other popular platforms that are sometimes referred to as social media services include YouTube, Letterboxd, QQ, Quora, Telegram, WhatsApp, Signal, LINE, Snapchat, Viber, Reddit, Discord, and TikTok. Wikis are examples of collaborative content creation.
Social media outlets differ from old media (e.g. newspapers, TV, and radio broadcasting) in many ways, including quality, reach, frequency, usability, relevancy, and permanence. Social media outlets operate in a dialogic transmission system (many sources to many receivers) while traditional media operate under a transmission model (one source to many receivers). For instance, a newspaper is delivered to many subscribers, and a radio station broadcasts the same programs to a city.
Social media has been criticized for a range of negative impacts on children and teenagers, including exposure to inappropriate content, exploitation by adults, sleep problems, attention problems, feelings of exclusion, and various mental health maladies. Social media has also received criticism as worsening political polarization and undermining democracy. Major news outlets often have strong controls in place to avoid and fix false claims, but social media's unique qualities bring viral content with little to no oversight. "Algorithms that track user engagement to prioritize what is shown tend to favor content that spurs negative emotions like anger and outrage. Overall, most online misinformation originates from a small minority of “superspreaders,” but social media amplifies their reach and influence."
History
Early computing
The PLATO system was launched in 1960 at the University of Illinois and subsequently commercially marketed by Control Data Corporation. It offered early forms of social media features with innovations such as | Technology | Internet | null |
690415 | https://en.wikipedia.org/wiki/Cairo%20Metro | Cairo Metro | The Cairo Metro (, lit. "Cairo Tunnel Metro" or ) is a rapid transit system in Greater Cairo, Egypt. It was the first of the three full-fledged metro systems in Africa and the first in the Middle East to be constructed. It was opened in 1987 as Line 1 from Helwan to Ramses Square with a length of . As of 2013, the metro carried nearly 4 million passengers per day. As of 15 May 2024, the Cairo Metro has 84 stations of which 5 are transfer stations, with a total length of . The system consists of three operational lines numbered 1 to 3.
The Cairo Metro is owned by the National Authority for Tunnels. The lines use standard gauge ().
Operations
The middle two cars (4th and 5th) of each train have been reserved for women since 1989 (the 5th car becomes mixed-use after 21:00). There are blue signs (pink on the first and second lines) at every station that signify the position of these cars. These cars are used as an option for women who do not wish to ride with men in the same car; however, women can still ride other cars freely. This policy was introduced to the protection of women from sexual harassment by men.
Cairo Metro operates from 05:00 to 01:00 (except during Ramadan, when it operates from 05:00 to 02:00), with the remaining hours reserved for maintenance work.
The ticket price was £E1 for each journey, regardless of distance. The Ministry of Transportation agreed to double the ticket prices, starting from Friday 24 March 2017, costing £E2 for the normal ticket, £E1.5 for the ticket of the minor, £E1 for the special needs ticket after formerly costing £E1, £E0.75 and £E0.5, respectively. Again on 10 May 2018, the ministry of transportation agreed to raise the ticket prices to be £E3 for 9 stops, £E5 for 16 stops and £E7 for more than 16 stops. As of July 2020, 9 stops cost £E5, 16 stops cost £E7, and rides exceeding 16 stops cost £E10.
Network
Line 1
Line 1 (blue) is the oldest line of the Cairo Metro and the first metro to open in Africa and the Middle East, with its first segment having opened in 1987. The line is long and serves 35 stations. This line carries trains with 3 units (9 train cars), which have a headway of 3.5 to 4 minutes, and a maximum speed of . The line can carry 60,000 passengers per hour in each direction.
Line 1 has a train driving simulator supplied by Transurb Technirail that won the international tender issued by Cairo Metro in 2011.
Line 2
Line 2 (red) is the second line of the Cairo Metro, first opening in 1996. The line is long, of which are in tunnels. It serves 20 stations, of which 12 are underground. It is mostly in the bored tunnel, with two exceptions: a short section at the northern end approaching Shubra El Kheima which is elevated, and a section just south of this by cut-and-cover. Line 2 uses the third rail electrification system instead of the overhead line used in the first line. The communication extension for line 2 was provided by Alcatel in 2005.
The minimum headway for the line is 2.7 to 3 minutes.
Line 2 has a simulator installed in Shubra since 2002 which was delivered by French company CORYS.
The first tunnel to be built under the Nile River carries Line 2 across the river.
Line 3
Line 3 (green) is the newest line. Its first section opened in 2012. It operates from Rod El-Farag Corridor and Cairo University to Adly Mansour, where it connects to the Cairo Light Rail Transit. Eventually, it will be extended to Cairo International Airport. The line crosses under the two branches of the River Nile, as does Line 2. The total length of the line will be approximately , most of which are in the bored tunnel, and will be implemented in four phases.
Phase 1 from Attaba station to Abbassia station opened in 21 February 2012, with five stations and a total length of . Phase 2 to Al Ahram Station was opened in 7 May 2014, by Adly Mansour, with four additional stations and an added length of , for a total length of . Phase 4A from Al Ahram to El Shams Club opened in 15 June 2019 (except for Heliopolis Square station which opened later that year). Phase 3A from Attaba station to Kit Kat station opened in October 2022. Phase 3B from Kit Kat station to Rod El-Farag Corridor station opened in January 2024, adding six more stations and . and began revenue service on 1 January 2024. Phase 3C from Kit Kat to Cairo University was opened in 15 May 2024.
History
Background
As the biggest and most densely populated megacity in Africa and the Middle East, Greater Cairo had a strong case for a metro. In 1987 that population stood at 10 million residents, not counting the two million or so commuters who came into Cairo every day to work. The capacity of Cairo's public transport infrastructure was around 20,000 passengers/hour, which increased to 60,000 after the construction of the metro.
Proposed plans
The idea of a metro was first proposed in the 1930s by engineer Saiyed Abdel Wahed of the Egyptian Railway Authority, however, the idea did not progress. Following the Egyptian Revolution of 1952, there was renewed interest in the idea. In 1954 French experts made a report about the future of the transportation in Egypt. They proposed a metro encompassing two lines, one long line connecting Bab al-Louq and Ismailia and a second line connecting Boulaq and Abou al-Ela Castle. They also proposed that there should be one company in charge of all transportation systems.
Later on, multiple experts came to Egypt regarding that project: Soviet experts in 1956, Japanese experts in 1960 and French experts in 1962, which concluded the following:
The creation of a metro system with multiple lines. The first is a line connecting Helwan with El-Marg going under Kasr el Eini street and Ramses Street. The second line would be from Sayyidah Zaynab Mosque to Shobra going under Downtown Cairo and would be long. The third would be from Giza to Abbaseya and would be long. The fourth would be from Al Awqaf to the Castle and would be long.
In 1964, British experts advised the creation of a metro line from Bab El Louk to Shubra. In 1966 Japanese experts advised the creation of a lane between Helwan and El Marg and another one going through Mohandiseen, Heliopolis and 26 of July Street, for a total of . Lastly, in 1969, the government approved the need for a study showing the needed capacity for Cairo's transportation system. The Egyptian Ministry of Transport issued an international tender for the creation of the study. Eight companies applied and the French company SOFRETU won the tender on 20 September 1970. The study was finished in 1973 and included mainly the study of the population development in Cairo and its needs for transportation capacity in 1980, 1985 and 1990. It concluded the necessity for three lines in greater Cairo to solve the transportation problem.
The first line would use the already available railways (Helwan-Bab El luk and Al Laymoun Bridge-El Marg) and connect them through a metro. It would be in total long. The second line would be and connect Shubra El-Kheima and Bulaq going through Ramses Street and Tahrir Square. The third line would go from El Darasa to Imbaba and would be in total long.
The priority for the project was the first line which would reduce 30% of the daily transport to and from Cairo and would incorporate pre-existing rail infrastructure. The detailed study of the construction took 6 years from 1975 to 1981. A tender was made for the construction of the metro and the Egyptian-French company Entra Nevra Arabco won it and had the task of constructing the metro system.
Construction works
The construction of Line 1 started in 1982 after the French government agreed on giving Egypt the necessary loan. The first section was opened in 27 September 1987 and the line was completed in 1989 connecting Helwan with El Marg and consisting of 34 stations with a total length of 42.5 km of which 4.7 km underground. In 1999, New El Marg station was added to the northern end of the line, bringing its total length to 44 km. Helwan University station was built between Wadi Houf and Ain Helwan stations.
Cairo's metro network was greatly expanded in the mid-1990s with the building of Line 2, from Shoubra El Kheima to Cairo University, with an extension to Giza. The line includes the first tunnel under the Nile. The construction of the line was finished in October 2000, and it was later extended to El Mounib.
Proposed lines
Line 4 (October-Oasis Highway - the Police Academy)
Line 4 is planned to run from Haram District to the New Cairo district, connecting Greater Cairo from West to East. It will cross the two branches of the Nile river and have a total length of .
Construction began in 2024 for completion in 2028.
Phase 1 (west) of the project will run from El-Malek El-Saleh Station (Interchange with Line 1) to the October-Oasis Highway Station with a total length of 18 km, passing through Giza Railway Station (Interchange with Line 2); the original plan for phase 1 was for it to start from El-Malek El-Saleh Station and end at the Grand Egyptian Museum Station with a total length of 10 km, but the Ministry of Roads & Transportation decided to extend the Line in their efforts to further connect the Governorate of 6 October to the Greater Cairo Area; Phase 1 also includes the plan to connect the end of Line 4 to the suburbs of 6 October mainly through executing The October 6th Tram system (The O6T) which will be by using a tram-train system supplied with the Alstom Regio-Citadis trams. This phase will have 15 stations to be constructed with a duration of 6.5 years. Phase 1 stations will be equipped with an automatic fare collection system and platform screen doors, and will include elevators for the use of disabled passengers.
Phase 1 bidding was postponed until May 2015 to enable Japan International Cooperation Agency (JICA) to complete the feasibility study and to resolve other problems with the construction starting by 2015 or 2016 according to Ismail El-Nagdy, Chairman of the National Authority for Tunnels with Japan International Cooperation Agency financing a $1.2 billion loan, while Egyptian government would cover the remainder, $2.4 billion.
Phase 2 will begin from El-Malek El-Saleh Station, passing east through Magra El-Oyoun street and Salah El-Din Citadel in Salah Salem street and ending at the 6th District Station (beginning of Nasr City district) with a completion date set at October 2018.
Phase 3 will begin from 6th District Station and end at Makram Ebeid Station, following Mustafa El-Nahas street in bored tunnels under the existing old tram system, and ultimately deconstructing the old railway and paving its right-of-way (increasing the street by two lanes in each direction, which is critically needed to lighten the traffic congestion in the area). This phase has a completion date set for October 2019. Mitsubishi was the only company to submit an offer on a tender for 64 trains in the third and fourth phases of the project. The offer was accepted in October 2018.
Phase 4, the final phase, will begin from Makram Ebeid Station with bored tunnels following Doctor Hassan El-Sherif street and Ahmed El-Zomor street, ending at Police Academy Station near the Ring Road.
The New Cairo Monorail
This project was first proposed by private investors as a plan to connect Line 3 with Line 4, through a route that is mostly parallel to the Ring Road's eastern arc, therefore covering New Cairo from north to south starting at the Cairo International Airport and ending at the beginning of the Cairo-Ain Sokhna Highway.
In August 2016, a different monorail project was under discussion, connecting 6 October city with other western parts of Greater Cairo. The line was expected to be long, cost $1.5 billion, and be completed by 2018.
A contract was signed in August 2019 with Bombardier Transportation, Orascom Construction, and Arab Contractors for $4.5 billion to build and operate the two monorail lines. The first will run 54 km from eastern Cairo to the new administrative capital. The second will run 42 km from 6 October City to Giza.
Qalyub Line 2 extension
As of December 2017, Spokesperson for the Egyptian National Authority for Tunnels (NAT) said that there are plans to extend Line 2 7 kilometers to the north from Shubra Al-Kheima station to end at Qalyub station due to the increasing traffic in north Cairo entrances, He also said that bidding will be held in February 2018 and the winner will be revealed mid-2018. The station will be an at-grade station built parallel to the train railway, bridges will be built to divert traffic from the Metro route.
Cairo International Airport Line 3 extension
Line 3 Phase 4C is proposed to start from Heliopolis station passing through Al-Hegaz Square and Military Academy area to Sheraton District and after ending at Cairo International Airport. If built, it would have five tunnel stations. There is no proposed completion date for this phase.
Long-term plans
A transportation study of the Greater Cairo region was completed in 1999. It recommended the implementation of a six-line system consisting of lines 1, 2 and 3 (existing); and lines 4, 5 and 6. The completed Metro Network would be capable of serving most of the densely populated areas in the Greater Cairo region, which was much in need of a comprehensive mass transit system. The plans include interchange stations between the six metro lines and would also provides interchange facilities with existing main railway stations, the airport, and bus stations.
The six planned metro lines aim to meet the transportation demands of the Greater Cairo area up to the year 2032. However, the actual construction and implementation schedule will be restricted by available funding, and the timetable will likely slip.
Line 5 would be a half-circular line connecting lines 1-4 in northern Cairo, running from Nasr City in the east to Port Said Street and Shubra El Kheima in the west. It would have a length of , entirely within bored tunnels. It would intersect Line 1 at Helmiet el-Zaitoun station, Line 2 at El-Khalafawy station, Line 3 at Haroun station, and Line 4 at Al-Wafaa we al-Amal station.
Line 6 would be a north-south line, from Shubra in the north to the Maadi and Helwan districts in the south. It would run from Ataba Station through El Kalaa street in bored tunnels to Salah Eldin Citadel Station (Interchange with Line 4) and move on from there to both districts via bored tunnels using the existing route El-Mahager Railway as a guide through both Maadi and Helwan. This Line has a length of . It would intersect Line 1 at Ghamra station, Line 3 at Bab al-Sharia station, Line 4 at Amr Ibn al-Ase station, and Line 5 at Sawah station.
Network map
| Technology | Africa | null |
691224 | https://en.wikipedia.org/wiki/Hot%20flash | Hot flash | Hot flashes, also known as hot flushes, are a form of flushing, often caused by the changing hormone levels that are characteristic of menopause. They are typically experienced as a feeling of intense heat with sweating and rapid heartbeat, and may typically last from two to 30 minutes for each occurrence.
Signs and symptoms
Hot flashes, a common symptom of menopause and perimenopause, are typically experienced as a feeling of intense heat with sweating and rapid heartbeat, and may typically last from two to thirty minutes for each occurrence, ending just as rapidly as they began. The sensation of heat usually begins in the face or chest, although it may appear elsewhere such as the back of the neck, and it can spread throughout the whole body. Some people feel as if they are going to faint. In addition to being an internal sensation, the surface of the skin, especially on the face, becomes hot to the touch. This is the origin of the alternative term "hot flush", since the sensation of heat is often accompanied by visible reddening of the face.
The hot-flash event may be repeated a few times each week or every few minutes throughout the day. Hot flashes may begin to appear several years before menopause starts and last for years afterwards. Some people undergoing menopause never have hot flashes. Others have mild or infrequent flashes. Those most affected experience dozens of hot flashes each day. In addition, hot flashes are often more frequent and more intense during hot weather or in an overheated room, the surrounding heat apparently making the hot flashes themselves both more likely to occur, and more severe.
Severe hot flashes can make it difficult to get a full night's sleep (often characterized as insomnia), which in turn can affect mood, impair concentration, and cause other physical problems. When hot flashes occur at night, they are called "night sweats". As estrogen is typically lowest at night, some people get night sweats without having any hot flashes during the daytime.
Young females
If hot flashes occur at other times in a young female's menstrual cycle, then it might be a symptom of a problem with the pituitary gland; seeing a doctor is highly recommended. In younger females who are surgically menopausal, hot flashes are generally more intense than in older females, and they may last until natural age at menopause.
Males
Hot flashes in males could have various causes. It can be a sign of low testosterone. Males with prostate cancer or testicular cancer can also have hot flashes, especially those who are undergoing hormone therapy with antiandrogens, also known as androgen antagonists, which reduce testosterone to castrate levels. Males who are castrated can also get hot flashes.
Types
Some menopausal females may experience both standard hot flashes and a second type sometimes referred to as "slow hot flashes" or "ember flashes". The standard hot flash comes on rapidly, sometimes reaching maximum intensity in as little as a minute. It lasts at full intensity for only a few minutes before gradually fading.
Slow "ember" flashes appear almost as quickly but are less intense and last for around half an hour. Females who experience them may undergo them year round, rather than primarily in the summer, and ember flashes may linger for years after the more intense hot flashes have passed.
Mechanism
Research on hot flashes is mostly focused on treatment options. The exact cause and pathogenesis, or causes, of vasomotor symptoms (VMS)—the clinical name for hot flashes—has not yet been fully studied.
Hot flashes are associated with declining levels of estrogen (estrogen withdrawal) and other hormonal changes. It does not appear that low levels of estrogen are the sole cause of hot flashes, as women who experience hot flashes have around the same plasma estrogen levels as women who do not have them, and prepubertal girls do not have hot flashes despite low estrogen levels.
There are indications that hot flashes may be due to a change in the hypothalamus's control of temperature regulation.
Transgender men also commonly report experiences of hot flashes. This is linked to hormonal changes possible from many aspects of masculinizing gender-affirming care, including the use of gonadotropin-releasing hormone agonists as puberty blockers, reduction of estrogen levels after having undergone oopherectomy, and long term testosterone use reducing production of estradiols.
Treatment
Hormone replacement therapy
Hormone replacement therapy may relieve many of the symptoms of menopause. However, oral HRT may increase the risk of breast cancer, stroke, and dementia and has other potentially serious short-term and long-term risks. Since the incidence of cardiovascular disease in women has shown a rise that matches the increase in the number of post menopausal women, recent studies have examined the benefits and side effects of oral versus transdermal application of different estrogens and found that transdermal applications of estradiol may give the vascular benefits lowering the incidences of cardiovascular events with fewer adverse side effects than oral preparations.
Women who experience troublesome hot flashes are advised by some to try alternatives to hormonal therapies as the first line of treatment. If a woman chooses hormones, they suggest she take the lowest dose that alleviates her symptoms for as short a time as possible. The US Endocrine Society concluded that women taking hormone replacement therapy for 5 years or more experienced overall benefits in their symptoms including relief of hot flashes and symptoms of urogenital atrophy and prevention of fractures and diabetes.
When estrogen as estradiol is used transdermally as a patch, gel, or pessary with micronized progesterone this may avoid the serious side effects associated with oral estradiol HRT since this avoids first pass metabolism (Phase I drug metabolism). Women taking bioidentical estrogen, orally or transdermally, who have a uterus must still take a progestin or micronized progesterone to lower the risk of endometrial cancer. A French study of 80,391 postmenopausal women followed for several years concluded that estrogen in combination with micronized progesterone is not associated with an increased risk of breast cancer. The natural, plant-derived progesterone creams sold over the counter contain too little progesterone to be effective. Wild yam (Dioscorea villosa) extract creams are not effective since the natural progesterone present in the extract is not bioavailable.
Selective serotonin reuptake inhibitors
SSRIs are a class of pharmaceuticals that are most commonly used in the treatment of depression. They have been found efficient in alleviating hot flashes. On 28 June 2013 FDA approved Brisdelle (low-dose paroxetine mesylate) for the treatment of moderate-to-severe vasomotor symptoms (e.g. hot flashes and night sweats) associated with menopause. Paroxetine became the first and only non-hormonal therapy for menopausal hot flashes approved by FDA.
Clonidine
Clonidine is a blood pressure-lowering medication that can be used to relieve menopausal hot flashes when hormone replacement therapy is not needed or not desired. For hot flashes, clonidine works by helping reduce the response of the blood vessels to stimuli that cause them to narrow and widen. While not all women respond to clonidine as a hot flash medication, it can reduce hot flashes by 40% in some peri-menopausal women.
Isoflavones
Isoflavones are commonly found in legumes such as soy and red clover. The two soy isoflavones implicated in relieving menopausal symptoms are genistein and daidzein, and are also known as phytoestrogens. The half life of these molecules is about eight hours, which might explain why some studies have not consistently shown effectiveness of soy products for menopausal symptoms. Although red clover (Trifolium pratense) contains isoflavones similar to soy, the effectiveness of this herb for menopausal symptoms at relatively low concentrations points to a different mechanism of action.
Other phytoestrogens
It is believed that dietary changes that include a higher consumption of phytoestrogens from sources such as soy, red clover, ginseng, and yam may relieve hot flashes.
Ginseng: Very few studies exist on the effect of ginseng for relief of menopausal symptoms. In a large double-blinded randomized controlled trial, reduction in hot flashes was not statistically significant but showed a strong trend towards improvement. Lack of statistical significance suggests future research, but does not meet the scientific bar for ginseng to be deemed effective.
Flaxseed: There have also been several clinical trials using flaxseed. Flaxseed is the richest source of lignans, which is one of three major classes of phytoestrogen. Lignans are thought to have estrogen agonist and antagonist effects as well as antioxidant properties. Flaxseed and its lignans may have potent anti-estrogenic effects on estrogen receptor positive breast cancer and may have benefits in breast cancer prevention efforts. One recent study done in France, looked at four types of lignans, including that found in flaxseed (Secoisolariciresinol) in a prospective cohort study to see if intake predicted breast cancer incidence. The authors report lowered risk of breast cancer among over 58,000 postmenopausal women who had the third highest quartile of lignan intake. There have been a few small pilot studies that have tested the effect of flaxseed on hot flashes. Currently there is a large study sponsored by the National Cancer Institute that is ongoing, but not accepting any new participants. The rationale for the study is that estrogen can relieve the symptoms of menopause, but can also cause the growth of breast cancer cells. Flaxseed may reduce the number of hot flashes and improve mood and quality of life in postmenopausal women not receiving estrogen therapy.
Acupuncture
Acupuncture has been suggested to reduce incidence of hot flashes in women with breast cancer and men with prostate cancer, but the quality of evidence is low.
Epidemiology
It has been speculated that hot flashes are less common among Asian women.
Menopausal women who have more hot flashes have a higher risk of hypertension and cardiovascular disease.
| Biology and health sciences | Human reproduction | Biology |
692463 | https://en.wikipedia.org/wiki/Rotation%20operator%20%28quantum%20mechanics%29 | Rotation operator (quantum mechanics) | This article concerns the rotation operator, as it appears in quantum mechanics.
Quantum mechanical rotations
With every physical rotation , we postulate a quantum mechanical rotation operator that is the rule that assigns to each vector in the space the vector
that is also in . We will show that, in terms of the generators of rotation,
where is the rotation axis, is angular momentum operator, and is the reduced Planck constant.
The translation operator
The rotation operator , with the first argument indicating the rotation axis and the second the rotation angle, can operate through the translation operator for infinitesimal rotations as explained below. This is why, it is first shown how the translation operator is acting on a particle at position x (the particle is then in the state according to Quantum Mechanics).
Translation of the particle at position to position :
Because a translation of 0 does not change the position of the particle, we have (with 1 meaning the identity operator, which does nothing):
Taylor development gives:
with
From that follows:
This is a differential equation with the solution
Additionally, suppose a Hamiltonian is independent of the position. Because the translation operator can be written in terms of , and , we know that This result means that linear momentum for the system is conserved.
In relation to the orbital angular momentum
Classically we have for the angular momentum This is the same in quantum mechanics considering and as operators. Classically, an infinitesimal rotation of the vector about the -axis to leaving unchanged can be expressed by the following infinitesimal translations (using Taylor approximation):
From that follows for states:
And consequently:
Using
from above with and Taylor expansion we get:
with the -component of the angular momentum according to the classical cross product.
To get a rotation for the angle , we construct the following differential equation using the condition :
Similar to the translation operator, if we are given a Hamiltonian which rotationally symmetric about the -axis, implies . This result means that angular momentum is conserved.
For the spin angular momentum about for example the -axis we just replace with (where is the Pauli Y matrix) and we get the spin rotation operator
Effect on the spin operator and quantum states
Operators can be represented by matrices. From linear algebra one knows that a certain matrix can be represented in another basis through the transformation
where is the basis transformation matrix. If the vectors respectively are the z-axis in one basis respectively another, they are perpendicular to the y-axis with a certain angle between them. The spin operator in the first basis can then be transformed into the spin operator of the other basis through the following transformation:
From standard quantum mechanics we have the known results and where and are the top spins in their corresponding bases. So we have:
Comparison with yields .
This means that if the state is rotated about the -axis by an angle , it becomes the state , a result that can be generalized to arbitrary axes.
| Physical sciences | Quantum mechanics | Physics |
692624 | https://en.wikipedia.org/wiki/Gotthard%20Base%20Tunnel | Gotthard Base Tunnel | The Gotthard Base Tunnel (GBT; , , ) is a railway tunnel through the Alps in Switzerland. It opened in June 2016 and full service began the following December. With a route length of , it is the world's longest railway and deepest traffic tunnel and the first flat, low-level route through the Alps. It lies at the heart of the Gotthard axis and constitutes the third tunnel connecting the cantons of Uri and Ticino, after the Gotthard Tunnel and the Gotthard Road Tunnel.
The GBT consists of a large complex with, at its core, two single-track tunnels connecting Erstfeld (Uri) with Bodio (Ticino) and passing below Sedrun (Grisons).
It is part of the New Railway Link through the Alps (NRLA) project, which also includes the Ceneri Base Tunnel further south (opened on 3 September 2020) and the Lötschberg Base Tunnel on the other main north–south axis.
It is referred to as a "base tunnel" since it bypasses most of the existing vertex line, the Gotthard railway line, a winding mountain route opened in 1882 across the Saint-Gotthard Massif, which was operating at its capacity before the opening of the GBT.
The new base tunnel establishes a direct route usable by high-speed rail and heavy freight trains.
The main purpose of the Gotthard Base Tunnel is to increase local transport capacity through the Alpine barrier, especially for freight on the Rotterdam–Basel–Genoa corridor, and more specifically to shift freight volumes from trucks to freight trains. This both significantly reduces the danger of fatal road crashes involving trucks, and reduces the environmental damage caused by heavy trucks. The tunnel also provides a faster connection between the canton of Ticino and the rest of Switzerland, as well as between northern and southern Europe, cutting the Basel/Zürich–Lugano–Milan journey time for passenger trains by one hour (and from Lucerne to Bellinzona by 45 minutes).
After 64 percent of Swiss voters accepted the NRLA project in a 1992 referendum, the first preparatory and exploratory work began in 1996. Construction began in November 1999 at Amsteg. Drilling operations were completed in March 2011. Completed in 2016, the final cost was reported to be CHF 12.2 billion ( billion).
Description
The Gotthard Base Tunnel, with a length of and a total of of tunnels, shafts and passages, is the longest railway tunnel in the world, with a geodetic distance of between the two portals. It is also the first flat route through the Alps or any other major mountain range, with a maximum elevation of above sea level, corresponding to that of Bern. It is the deepest railway tunnel in the world, with a maximum depth of , comparable to that of the deepest mines on Earth. Without ventilation, the temperature inside the mountain reaches .
Like the two other tunnels passing below the Gotthard, the Gotthard Base Tunnel connects two Alpine valleys across the Saint-Gotthard Massif: the Urner Reusstal in the canton of Uri, in which flows the river Reuss, and the Valle Leventina, the largest valley in the canton of Ticino, in which the river Ticino flows. Unlike most other tunnels, the Gotthard Base Tunnel passes under several distinct mountain massifs, two of them being major subranges of the Alps, the Glarus Alps and the Saint-Gotthard Massif, with the valley of the Anterior Rhine, the Surselva in the canton of Graubünden, between them. The tunnel passes under these two ranges more than below the Chrüzlistock () and the Piz Vatgira (, near the Lukmanier Pass). While the cantons of Uri and Ticino are part of the German- and Italian-speaking areas of Switzerland respectively, the Surselva is mainly Romansh-speaking.
The Alps strongly influence the European climate – and that of Switzerland in particular – and there can be substantially different weather conditions at each end of the GBT, described by the Ticinese architect Mario Botta: "The light changes at the Gotthard: that of the Mediterranean Sea is not the same as that of the continent, that of the central lands, that of Europe far away from the sea." On average, the temperature is about higher on the south side than the north side, but on some days, temperature differences are well over .
The north portal lies in the north of the municipality of Erstfeld at an elevation of , east of the Reuss. There, the tunnel penetrates the western slopes of the Bälmeten and Chli Windgällen (although only marginally) before passing below the valley of the Chärstelenbach, a creek in the Maderanertal. From there, the tunnel runs parallel to the small valley of Etzli, below the Witenalpstock. The main crest of the Glarus Alps, which is the watershed between the Reuss and the Anterior Rhine, is crossed below the Chrüzlistock, the crest having an elevation of about at this point. From the crest and border, the tunnel runs parallel to the small valley of the river Strem (Val Strem) before passing below Sedrun and the Anterior Rhine. From the bottom of the valley, the tunnel proceeds towards the valley of the Rein da Nalps (Val Nalps) and passes east of Lai da Nalps, before crossing the Gannaretsch range below the western summit of Piz Vatgira (). This is the deepest point of the tunnel, with a rock layer of above it. The tunnel then passes below the valley of the Rein da Medel (Val Medel) and west of Lai da Sontga Maria. After a few kilometres the tunnel crosses the watershed between the Anterior Rhine and the Ticino, just north of Pizzo dell'Uomo (). This point corresponds to the main chain of the Alps, and is the main drainage divide between the Rhine and the Po. For a few kilometres, the tunnel passes below two western tributaries of the Brenno in the Valle Santa Maria before crossing the last range, west of the Passo Predèlp (about ) and east of Faido. It then follows the eastern slopes of the large Valle Leventina, the valley of the Ticino, for about to the south portal at Bodio, at an elevation of , just before Biasca, where the Brenno converges with the Ticino.
The closest railway stations to the portals are Altdorf and Biasca. The first regularly served railway stations on the base line (as of 2016/17) are those of Arth-Goldau (Schwyz), a railway node with links to Lucerne and Zürich, and Bellinzona (the "Gate of Ticino"), with links to Locarno, Luino and Lugano (via the Monte Ceneri Rail Tunnel). The journey from Arth-Goldau to Bellinzona takes not more than an hour. The station of Altdorf is planned to be served by 2021. There also have been talks of using that of Biasca. The travel between Altdorf and Biasca would last less than 25 minutes.
History
Background
Since the 13th century, the Gotthard Pass has been an important trade route from northern to southern Europe. Control of its access routes led to the birth of the Swiss Confederacy. The Gotthard Pass is located halfway between Lake Lucerne and Lake Maggiore. It is the shortest link between the navigable Rhine and the Po. Before modern transport, the traverse of the pass took days, and snow makes it a challenge in winter.
Quite late, compared to other current top-importance routes through the Alps (e.g. Simplon, San Bernardino, Brenner, Mont Cenis), namely in 1830, the first Saint-Gotthard Pass road was established after centuries-long usage of a bridle path. From 1842 onwards, a daily course by the Gotthard Post, a stagecoach drawn by five horses with ten seats, still took about 23 hours from Como to Flüelen. It would last until 1921.
In 1882, with the inauguration of the Gotthard Railway Tunnel, the travel time between Altdorf and Biasca was reduced dramatically to only hours, though often accompanied with overnight stays in large Fin de siècle-hotels, for example in Biasca. In those days, it was still an adventure and it was only affordable to the rich. Electrification of the railway line in 1922 significantly reduced travel time again. Refilling water boilers of steam locomotives was no longer necessary. There were also the technical advantages of electrical engines and future technical improvements.
From 1924, car transport on trains through the railway tunnel began. The road between Göschenen and Airolo over the summit of the pass, comporting notably the Schöllenen ravine and the Tremola, had countless hairpin turns and serpentine curves, dropping in altitude. It posed a huge challenge for automobiles of those days. From 1953 onwards, the pass road was sequentially improved and expanded at several sections along the Gotthard route, finally ending in 1977 with the opening of an expressway fully circumventing the Tremola. In winter, however, due to the snow, cars could only cross the Gotthard on the train.
Transit time was further dramatically reduced with the opening of the Gotthard Road Tunnel and the finalization of the northern part of A2 motorway through the Urner Reusstal (in close proximity to the railway), with many additional tunnels (then leading from Basel to the Gotthard Road Tunnel), in 1980. With the completion in 1986 of the A2 motorway in the Valle Leventina, the main valley leading from Airolo down to Bellinzona, and the surmounting of the Monte Ceneri between Bellinzona and Lugano in 1983, finally a continuous motorway was established from the northern border of Switzerland in Basel to the southern border in Chiasso, or the shortest motorway route from North-German Hamburg as far as South-Italian Sicily, bringing down the competitiveness of the railway line. Passenger speed was also increased on the railway line with the use of tilting trains, notably the ICN, although maximum speed remaining far lower than on a modern straight high-speed line. Both modern motorway and historic railway rely on heavy rockfall and avalanche protection equipments and are exposed to harsh weather condition in winter.
After the opening of the auto tunnel, in 1980, traffic increased more than tenfold. The existing tunnel was at its capacity by 2013. A second tunnel will be built next to the first, following a national referendum. Construction started in 2021 and is scheduled to finish in 2027.
As early as 1947, engineer Eduard Gruner imagined a two-story base tunnel from Amsteg to Biasca, both rail and road, with a stop at Sedrun, to provide a faster and flatter passage through the Swiss Alps. Similarly to Gruner's idea, the GBT cuts through the Gotthard Massif some below the older tunnel. On the historic track only the Gotthard Railway trains up to when using two locomotives or up to with an additional bank engine at the end of the train are able to pass through the narrow mountain valleys and through spiral tunnels climbing up to the portals of the old tunnel at a height of above sea level. Since the GBT is in full service, standard freight trains of up to are able to pass this natural barrier.
Because of ever-increasing international truck traffic, Swiss voters chose a shift in transportation policy in September 1992 by accepting the NRLA proposal. A second law, the Alpine Protection Act of February 1994, requires a shift of as much tonnage as possible from truck transport to train transport.
The goal of both the laws is to transport trucks, trailers and freight containers through Switzerland, from Basel to Chiasso, and beyond by rail to relieve the overused roads, and that of the Gotthard in particular, by using intermodal freight transport and rolling highways (where the entire truck is transported). The GBT substantially contributes to the requirements of both laws and enables a direct flat route from the ports of the North Sea (notably Rotterdam) to those of the Mediterranean Sea (notably Genoa), via the Rhine corridor.
Although the technical maximum speed is through the GBT, the maximal authorized speed has been reduced to for ecological and economical reasons, while the operating speed of passenger trains is restricted to in order to accommodate the freight traffic, with the possibility to accelerate up to in case of delay. At opening the GBT reduced travel times for trans-Alpine train journeys by about 40 minutes, and by one hour when the adjacent Zimmerberg and Ceneri Base Tunnels were completed. This is viewed as a revolution, especially in the isolated region of Ticino, which is separated from the rest of the country by the Alps and the Gotthard. The two stations of Bellinzona and Lugano (respectively named "Gate of Ticino" and "Terrace of Ticino") were entirely renovated for the opening of the GBT, among other improvements.
As of 2016, the Gotthard Base Tunnel is the longest railway tunnel in the world. It is the third Swiss tunnel to bear this title, after the Gotthard Tunnel (, 1882) and the Simplon Tunnel (, 1905). It is the third tunnel built under the Gotthard, after the Gotthard Tunnel and the Gotthard Road Tunnel.
Construction
AlpTransit Gotthard AG was responsible for construction. It is a wholly owned subsidiary of the Swiss Federal Railways (SBB CFF FFS).
To cut construction time in half, four access tunnels were built so that construction could start at four different sites simultaneously: Erstfeld, Amsteg, Sedrun, and Faido. A fifth at Bodio was added later. The two tunnels are joined approximately every by connecting galleries. Trains can move between the tunnels in the two multifunction stations at Sedrun and Faido. These stations house ventilation equipment and technical infrastructure and serve as emergency stops and evacuation routes.
Access to the Sedrun station site is by a level access tunnel long from the valley floor near Sedrun. At the end of the access tunnel, two vertical shafts lead down to the base tunnel level. A proposal to construct a functioning railway station, called Porta Alpina (from Romansh, "Alpine Gate"), at this site was evaluated, but the project was put on hold in 2007 and definitively cancelled by the federal authorities in 2012 as uneconomical.
The final breakthrough in the east tube occurred on 15 October 2010 at 14:17 +02:00. The final breakthrough in the west tube occurred on 23 March 2011 at 12:20.
On 30 August 2013, the tunnel was entirely traversed for the first time from Bodio to Erstfeld in six hours, by diesel train, buses and by foot.
On 16 December 2013, the operational test phase started on a stretch in the southern section of the west tube between Faido and Bodio. Its purpose was to test the infrastructure and any ancillary systems.
On 31 October 2014, the railway track installation was completed. A gold sleeper on the very last part of the track was installed during the event to mark this milestone of progress.
On 1 October 2015, following permission by the Federal Office of Transport, the first tests on the entire length of the GBT were performed, with steadily increasing speed. On 8 November, a train reached the top speed of .
Allocation of work
The contracts were awarded in sections:
Erstfeld (the section from Erstfeld to Amsteg), with two tunnel boring machines (TBM) boring the two tubes. The break-through of the east tube between Erstfeld and Amsteg took place on 15 June 2009. The portal area was surface-mined.
Amsteg (the section from Amsteg to north of Sedrun), ARGE AGN (Strabag and Züblin Murer) received the contract for work in this sector. On 9 December 2009, the Amsteg section was officially delivered to the owner for fitting-out, with civil engineering, construction, concrete and lining work completed in early 2010.
Sedrun (the East tube and West tube in the section immediately north and south of Sedrun), along with work performed by Transco (Bilfinger SE, Implenia, Frutiger and Impresa Pizzarotti). The final breakthrough in the west tube occurred in March 2011. The northbound tubes from Amsteg to the Sedrun multifunction station (north) were handed over to the railway systems contractor Transtec Gotthard on 15 September 2011, the date specified in the construction schedule.
Faido ( East tube and West tube in the section from south of Sedrun to Faido), with Consorzio TAT (Alpine Mayreder Bau, CSC Impresa costruzioni, Hochtief and Implenia and Impregilo).
Bodio ( East tube and West tube in the section from Faido to Bodio), with work performed by Consorzio TAT (Alpine Mayreder Bau, CSC Impresa costruzioni, Hochtief, Implenia and Impregilo). Civil engineering construction, concrete and lining works were completed in early 2010.
Deaths during construction
Nine workers died during construction; one in the Amsteg section, two in the Sedrun section, and three each in the southernmost Faido and Bodio sections.
Inauguration and commissioning
In 2016, several events, including festivities and special exhibitions, were held around the Gotthard, culminating in the inaugurations in early June, dubbed Gottardo 2016. Public institutions joined the celebrations: Swiss Post issued a special stamp commemorating the Gotthard Base Tunnel, and Swissmint issued gold and silver coins dedicated to the opening.
On 31 May 2016, a day before the inauguration, the nine people who died during construction were commemorated in a ceremony at the north portal in Erstfeld that was led by a Catholic vicar general, a vicar of the Evangelical-Reformed Church of Uri, a Jewish rabbi, and a Muslim imam. A bronze memorial plaque with their names — four coming from Germany, three from Italy, and one from each of South Africa and Austria – was unveiled by AlpTransit Gotthard AG CEO Renzo Simoni. A Catholic shrine to Saint Barbara, the patron of miners, stands inside the tunnel as a memorial.
The tunnel was officially inaugurated on 1 June 2016. At the northern entrance in Erstfeld, President of the Confederation Johann Schneider-Ammann spoke of a "giant step for Switzerland but equally for our neighbours and the rest of the continent", while a live relay carried a speech given by Transport Minister Doris Leuthard at the southern entrance in Bodio. The first journey carried hundreds of Swiss citizens who had won tickets in a draw, while the assembled guests in Erstfeld, including the Federal Council in corpore, heads of state and government from neighbouring countries and transport ministers from European countries, attended the opening show Sacre del Gottardo by Volker Hesse featuring 600 dancers, acrobats, singers and musicians celebrating Alpine culture and myths around the Gotthard. On the following weekend, popular festivities and special exhibitions, attended by more than 100,000 visitors, were held.
From 2 August to 27 November 2016, the Swiss Federal Railways ran a special train service through the tunnel called "Gottardino" which was open to the public. It was a once-daily service from Flüelen railway station to Biasca railway station and in reverse. The trains made a stop inside the tunnel, to allow passengers to visit an exhibition inside the underground multifunction station in Sedrun which would normally be used in emergency only.
Regular services
During 2016, the GBT was tested extensively before its integration into the regular schedule on 11 December. On 5 December, the Swiss Federal Railways were granted permission from the Federal Transport Office to use the new base line. While the base tunnel is used for InterCity trains (ICN) and EuroCity trains (EC), the vertex line remains in use for regional trains. Since 2019, the Gotthard axis is served by the Stadler EC250 (Giruno) high-speed train and future flagship of the SBB fleet.
From the Amsteg portal, guided tours are organised inside the Gotthard Base Tunnel complex. A window allows visitors to watch the trains running in the tunnel.
2023 derailment and closure
On 10 August 2023, a freight train derailed while traveling through the tunnel, causing extensive damage to the tunnel infrastructure. The incident occurred near the multi-functional station of Faido, in the canton of Ticino. No one was injured, and no hazardous materials were released.
The derailment caused significant damage to the tunnel, prompting its closure to both passenger and freight traffic. The tunnel returned to normal service on 2 September 2024.
Politics
The realization of the GBT, as the centrepiece of the NRLA, is also a prototypical example of direct democracy in Switzerland. In order to accomplish this mega-project the political institutions also had to overcome many parliamentary sessions and several major popular votes, including the following:
27 September 1992, NRLA proposal (mandatory referendum): The final proposal by the Federal Council was accepted by 63.6% yes votes (declined by 1+2/2 cantons, turnout 45.9%)
20 February 1994, Alps Initiative (federal popular initiative): Initiated by a few private people with the goal to protect the Alpine environment from the negative impact of traffic was accepted by 51.9% yes votes (declined by 7 cantons, turnout 41%). The initiative was accepted despite the recommendation by the Federal Council from 12 February 1992 to decline the initiative without any counterproposal, and despite the parliamentary recommendation (both chambers) from 18 June 1993 to decline the initiative.
29 November 1998, Public Transport Funding (mandatory referendum): A total budget of CHF 30 billion for several public transport projects was accepted by 63.5% yes votes (declined by 1+3/2 cantons, turnout 38.3%); "the NRLA is to receive CHF 13.6 billion"
21 May 2000, Bilateral EU Agreements / 40-tonne Trucks / Heavy Traffic Fee (optional referendum): As part of a whole package of several bilateral agreements with the EU the Swiss also accepted by 67.2% yes votes (declined by 2 cantons, turnout 48.3%) the shift of an upper limit for trucks from to , but at the same time the EU agreed to a new heavy-traffic fee, which would also be used to finance the NRLA
17 December 2003, Ceneri Base Tunnel (parliamentary session): The controversial funding of the Ceneri Base Tunnel was finally passed by parliamentary approval only; the possibility for an optional referendum was not raised by any political groups, nor by the public. The then-in-charge transport minister, Federal Councilor Moritz Leuenberger, was quoted as saying "This is the only way to make the railway [the Gotthard axis] a flat line between Basel and Chiasso."
Figures
Diameter of each of the single-track tubes:
Distance between cross passage tubes: ca.
Numbers of cross passage tubes: 178
Maximum rock overlay: (at Piz Vatgira)
Start of construction: 1993 (sounding drills), 1996 (preparations), 4 November 1999 (official start, first blasting), 2003 (mechanical excavation)
Breakthrough: 15 October 2010 (Eastern tube), 23 March 2011 (Western tube)
Commissioning: May 2016
Inauguration/opening: 1 June 2016
Start of daily passenger service: 11 December 2016 (see public transport timetable#Switzerland)
Total cost: CHF 9.560 billion ()
Travel time: Passenger trains – 20 minutes
Amount of excavated rock: , , the equivalent of 5 Giza pyramids
Number of tunnel boring machines (TBM): Four Herrenknecht Gripper TBMs. Machine numbers S-210 and S-211 operated northbound from Bodio to Faido and Sedrun and were nicknamed Sissi and Heidi respectively; Machines S-229 and S-230 operated southbound from Erstfeld to Sedrun and were known as Gabi I and Gabi II.
Total length: (including back-up equipment)
Total weight:
Power: 5 MW
Max. excavation daily: (in excellent rock conditions)
Total excavation length by TBM: about (for each tube)
Manufacturer: Herrenknecht, Schwanau, Germany
Operation
Reduced travel times
Safety
The safety requirements on the rolling stock are similar to those of other long Swiss tunnels, including the ability for the emergency brake to be overridden.
2023 freight train derailment
On 10th August 2023, a freight train heading north derailed inside the tunnel causing it to be shut. Trains scheduled to use the tunnel were either cancelled or redirected via the "panorama route" which added around one to two hours to journeys. Since the alternate route was unable to accommodate bi-level equipment, passenger capacity was reduced by around two thirds. According to an SBB press conference, around 8km of tracks, 20,000 concrete sleepers and a lane change gate in the tunnel were severely damaged. The gate is needed to separate the two tunnel tubes which is why the second tunnel could not initially be used to run trains.
Experts from the Swiss Transportation Safety Investigation Board (STSB) were dispatched to the site to preserve and document evidence of the crash. Investigators found fragments of a wheel and signs of derailment several kilometers before the crash site. At the Faido station, they discovered a pile of derailed wagons. The STSB's mission is to determine the causes of accidents to ensure or improve safety, rather than to assign guilt or responsibility.
The STSB believes that the derailment was likely caused by a broken wheel tread inside the tunnel. All fragments of the wheel were found, consisting of several large pieces. The cause of the breakage has not yet been determined. External influence or fatigue fracture are possible explanations. The wheel fragments will be subjected to metallurgical analysis to determine the cause.
Despite the broken wheel, the wagon was dragged by the train for several kilometers. At the switching point of the multi-functional station, where trains can switch between the two tubes of the tunnel, the wagon derailed, pulling more than 20 following wagons off the tracks.
The tunnel's repair work was extensive. Originally expected to be completed by the end of 2023, it was announced in November 2023 that the repairs would not be finished before September 2024, with the tunnel operating at reduced capacity until then. This was due to repairs requiring a complete replacement of seven kilometers of rails, including 20,000 sleeper blocks and their concrete foundations, a damaged lane-change gate, two high-speed switches and many other safety and operationally relevant parts. Freight trains were allowed to use the east tunnel from 23 August due to a temporary maintenance gate replacing the damaged one whilst passenger services remained rerouted.
Freight capacity in the tunnel was gradually increased on weekdays, while passenger services saw expansion on weekends. In case of emergency, passengers would be evacuated through the other tunnel tube. Following the reopening of the undamaged northern third of the western tube, additional freight and passenger services were scheduled to use the tunnel, with a full reopening and a return to normal service commencing on 2nd September 2024.
Traffic
Since the opening date on 1 June 2016, between 130 and 160 trains on an average working day operated through the Gotthard Base Tunnel, which in March 2019 marked the 100,000th transit. Around two-thirds of the passages were freight trains and the remaining quota were passenger trains, both national and international.
Projections
The number of projected trains per day was 180–260 freight trains and 50 (65 from 2020) passenger trains.
Passengers
After the opening of the tunnel there was an increase in passengers crossing the trans-alpine line, with 2.3 million passengers in the first 8 months, an increase of 30% over the previous year.
In August 2017, an average of 10,400 people crossed the tunnel daily. Train services from Italy to Switzerland through the line are expected to become faster from 2020, with the opening of the Ceneri Base Tunnel, with an expected further increase in passenger numbers.
There are plans for a train service between Zürich and Milan with a journey time of 2:45 hours, down from 3:50 hours.
Freight
on 120 trains passed through the tunnel each day during the first half year of operation.
| Technology | Tunnels | null |
692997 | https://en.wikipedia.org/wiki/Tannic%20acid | Tannic acid | Tannic acid is a specific form of tannin, a type of polyphenol. Its weak acidity (pKa around 6) is due to the numerous phenol groups in the structure. The chemical formula for commercial tannic acid is often given as C76H52O46, which corresponds with decagalloyl glucose, but in fact it is a mixture of polygalloyl glucoses or polygalloyl quinic acid esters with the number of galloyl moieties per molecule ranging from 2 up to 12 depending on the plant source used to extract the tannic acid. Commercial tannic acid is usually extracted from any of the following plant parts: Tara pods (Caesalpinia spinosa), gallnuts from Rhus semialata or Quercus infectoria or Sicilian sumac leaves (Rhus coriaria).
According to the definitions provided in external references such as international pharmacopoeia, Food Chemicals Codex and FAO-WHO tannic acid monograph only tannins obtained from the above-mentioned plants can be considered as tannic acid. Sometimes extracts from chestnut or oak wood are also described as tannic acid but this is an incorrect use of the term. It is a yellow to light brown amorphous powder.
While tannic acid is a specific type of tannin (plant polyphenol), the two terms are sometimes (incorrectly) used interchangeably. The long-standing misuse of the terms, and its inclusion in scholarly articles has compounded the confusion. This is particularly widespread in relation to green tea and black tea, both of which contain many different types of tannins not just exclusively tannic acid.
Tannic acid is not an appropriate standard for any type of tannin analysis because of its poorly defined composition.
Quercitannic and gallotannic acids
Quercitannic acid is one of the two forms of tannic acid found in oak bark and leaves. The other form is called gallotannic acid and is found in oak galls.
The quercitannic acid molecule is also present in quercitron, a yellow dye obtained from the bark of the Eastern black oak (Quercus velutina), a forest tree indigenous in North America. It is described as a yellowish-brown amorphous powder.
In 1838, Jöns Jacob Berzelius wrote that quercitannate is used to dissolve morphine.
In 1865 in the fifth volume of "A dictionary of chemistry", Henry Watts wrote :
It exhibits with ferric salts the same reactions as gallotannic acid. It differs however from the latter in not being convertible into gallic acid, and not yielding pyrogallic acid by dry distillation. It is precipitated by sulfuric acid in red flocks. (Stenhouse, Ann. Ch. Pharm. xlv. 16.)
According to Rochleder (ibid lxiii. 202), the tannic acid of black tea is the same as that of oak-bark.
In 1880, Etti gave for it the molecular formula C17H16O9. He described it as an unstable substance, having a tendency to give off water to form anhydrides (called phlobaphenes), one of which is called oak-red (C34H30O17). For him, it was not a glycoside.
In Allen's "Commercial Organic Analysis", published in 1912, the formula given was C19H16O10.
Other authors gave other molecular formulas like C28H26O15, while another formula found is C28H24O11.
According to Lowe, two forms of the principle exist – "one soluble in water, of the formula C28H28O14, and the other scarcely soluble, C28H24O12. Both are changed by the loss of water into oak red, C28H22O11."
Quercitannic acid was for a time a standard used to assess the phenolic content in spices, given as quercitannic acid equivalent.
In an interesting historical note, the inventor of carborundum, Edward G. Acheson, discovered that gallotannic acid greatly improved the plasticity of clay. In his report of this discovery in 1904 he noted that the only known historical reference to the use of organic material added to clay is the use of straw mixed with clay described in the Bible, Exodus 1:11 and that the Egyptians must have been aware of his (re-)discovery. He stated "This explains why the straw was used and why the children of Israel were successful in substituting stubble for straw, a course that would hardly be possible, were the fibre of the straw depended upon as a bond feasible for the clay, but quite reasonable where the extract of the plant was used."
Uses
Tannins are a basic ingredient in the chemical staining of wood, and are already present in woods like oak, walnut, and mahogany. Tannic acid can be applied to woods low in tannin so chemical stains that require tannin content will react. The presence of tannins in the bark of redwood (Sequoia) is a strong natural defense against wildfire, decomposition and infestation by certain insects such as termites. It is found in the seeds, bark, cones, and heartwood.
Tannic acid is a common mordant used in the dyeing process for cellulose fibers such as cotton, often combined with alum and/or iron. The tannin mordant should be done first as metal mordants combine well with the fiber-tannin complex. However this use has lost considerable interest.
Similarly tannic acid can also be used as an aftertreatment to improve wash fastness properties of acid dyed polyamide. It is also an alternative for fluorocarbon aftertreatments to impart anti-staining properties to polyamide yarn or carpets.
However, due to economic considerations currently the only widespread use as textile auxiliary is the use as an agent to improve chlorine fastness, i.e. resistance against dye bleaching due to cleaning with hypochlorite solutions in high-end polyamide 6,6-based carpets and swimwear. It is, however, used in relatively small quantities for the activation of upholstery flock; this serves as an anti-static treatment.
Tannic acid is used in the conservation of ferrous (iron based) metal objects to passivate and inhibit corrosion. Tannic acid reacts with the corrosion products to form a more stable compound, thus preventing further corrosion from taking place. After treatment the tannic acid residue is generally left on the object so that if moisture reaches the surface the tannic acid will be rehydrated and prevent or slow any corrosion. Tannic acid treatment for conservation is very effective and widely used but it does have a significant visual effect on the object, turning the corrosion products black and any exposed metal dark blue. It should also be used with care on objects with copper alloy components as the tannic acid can have a slight etching effect on these metals.
Tannic acid is also found in commercially available iron/steel corrosion treatments, such as Hammerite Kurust.
Use in food
In many parts of the world, its uses in food are permitted. In the United States, tannic acid is generally recognized as safe by the Food and Drug Administration for use in baked goods and baking mixes, alcoholic and non-alcoholic beverages, frozen dairy products, soft and hard candy, meat products, and rendered animal fat.
According to EU directive 89/107/EEC, tannic acid cannot be considered as a food additive and consequently does not hold an E number. Under directive 89/107/EEC, tannic acid can be referred to as a food ingredient. The E-number E181 is sometimes incorrectly used to refer to tannic acid; this in fact refers to the INS number assigned to tannic acid under the FAO-WHO Codex Alimentarius system.
Uses as a medication
In conjunction with magnesium and sometimes activated charcoal, tannic acid was once used as a treatment for many toxic substances, such as strychnine, mushroom, and ptomaine poisonings in the late 19th and early 20th centuries.
The introduction of tannic acid treatment of severe burn injuries in the 1920s significantly reduced mortality rates. During World War I, tannic acid dressings were prescribed to treat "burns, whether caused by incendiary bombs, mustard gas, or lewisite". After the war this use was abandoned due to the development of more modern treatment regimens.
Hazards
Tannic acid could cause potential health hazards such as damage to the eye, skin, respiratory tract, and gastrointestinal tract. It may cause irritation, redness, pain, blurred vision, and possible eye damage. When tannic acid is absorbed through the skin in harmful amounts, it may cause irritation, redness, and pain. Nausea, vomiting and diarrhoea are symptoms of tannic acid ingestion and prolonged exposure may cause liver damage. Upon inhalation, tannic acid may cause respiratory tract irritation.
Crocodilian coloration
Skin color in Crocodilia (crocodiles and alligators) is very dependent on water quality. Algae-laden waters produce greener skin, while tannic acid in the water from decay of leaves from overhanging trees (which produces some types of blackwater rivers) often produce darker skin in these animals.
| Physical sciences | Polyphenols | Chemistry |
693028 | https://en.wikipedia.org/wiki/Channichthyidae | Channichthyidae | The crocodile icefish or white-blooded fish comprise a family (Channichthyidae) of notothenioid fish found in the Southern Ocean around Antarctica. They are the only known vertebrates to lack hemoglobin in their blood as adults. Icefish populations are known to reside in the Atlantic and Indian sectors of the Southern Ocean, as well as the continental shelf waters surrounding Antarctica. Water temperatures in these regions remain relatively stable, generally ranging from . One icefish, Champsocephalus esox, is distributed north of the Antarctic Polar Frontal Zone. At least 16 species of crocodile icefish are currently recognized, although eight additional species have been proposed for the icefish genus Channichthys.
In February 2021, scientists discovered and documented a breeding colony of Neopagetopsis ionah icefish estimated to have 60 million active nests across an area of approximately 92 square miles at the bottom of the Weddell Sea in Antarctica. The majority of nests were occupied by one adult fish guarding an average of 1,735 eggs in each nest.
Genera
The following genera have been classified within the family Channichthyidae:
Chaenocephalus Richardson, 1844
Chaenodraco Regan, 1914
Champsocephalus Gill, 1861
Channichthys Richardson, 1844
Chionobathyscus Andriashev & Neyelov, 1978
Chionodraco Lönnberg, 1905
Cryodraco Dollo, 1900
Dacodraco Waite, 1916
Neopagetopsis , 1947
Pagetopsis Regan, 1913
Pseudochaenichthys Norman, 1937
Diet and body size
All icefish are believed to be piscivorous, but can also feed on krill. Icefish are typically ambush predators; thus, they can survive long periods between feeding, and often consume fish up to 50% of their own body length. Maximum body lengths of have been recorded in these species.
Respiratory and circulatory system
Icefish blood is colorless because it lacks hemoglobin, the oxygen-binding protein in blood. Channichthyidae are the only known vertebrates to lack hemoglobin as adults. Although they do not manufacture hemoglobin, remnants of hemoglobin genes can be found in their genome. The hemoglobin protein is made of two subunits (alpha and beta). In 15 of the 16 icefish species, the beta subunit gene has been completely deleted and the alpha subunit gene has been partially deleted. One icefish species, Neopagetopsis ionah, has a more complete, but still nonfunctional, hemoglobin gene.
Red blood cells (RBCs) are usually absent, and if present, are rare and defunct. Oxygen is dissolved in the plasma and transported throughout the body without the hemoglobin protein. The fish can live without hemoglobin via low metabolic rates and the high solubility of oxygen in water at the low temperatures of their environment (the solubility of a gas tends to increase as temperature decreases). However, the oxygen-carrying capacity of icefish blood is less than 10% that of their relatives with hemoglobin.
Myoglobin, the oxygen-binding protein used in muscles, is absent from all icefish skeletal muscles. In 10 species, myoglobin is found in the heart muscle, specifically ventricles. Loss of myoglobin gene expression in icefish heart ventricles has occurred at least four separate times.
To compensate for the absence of hemoglobin, icefish have larger blood vessels (including capillaries), greater blood volumes (four times those of other fish), larger hearts, and greater cardiac outputs (five times greater) compared to other fish. Their hearts lack coronary arteries, and the ventricle muscles are very spongy, which enables them to absorb oxygen directly from the blood they pump. Their hearts, large blood vessels and low-viscosity (RBC-free) blood are specialized to enable very high flow rates at low pressures. This helps to reduce the problems caused by the lack of hemoglobin. In the past, their scaleless skin had been widely thought to help absorb oxygen. However, current analysis has shown that the amount of oxygen absorbed by the skin is much less than that absorbed through the gills. The little extra oxygen absorbed by the skin may play a part in supplementing the oxygen supply to the heart, which receives venous blood from the skin and body before pumping it to the gills. Additionally, icefish have larger cardiac mitochondria and increased mitochondrial biogenesis in comparison to red-blooded notothenioids. This adaptation facilitates enhanced oxygen delivery by increasing mitochondrial surface area, and reducing distance between the extracellular area and the mitochondria.
Evolution
The icefish are considered a monophyletic group and likely descended from a sluggish demersal ancestor. The cold, well-mixed, oxygen-rich waters of the Southern Ocean provided an environment where a fish with a low metabolic rate could survive even without hemoglobin, albeit less efficiently.
When the icefish evolved is unknown; two main competing hypotheses have been postulated. The first is that they are only about 6 million years old, appearing after the Southern Ocean cooled significantly. The second suggests that they are much older, 15–20 million years.
Although the evolution of icefish is still disputed, the formation of the Antarctic Polar Frontal Zone (APFZ) and the Antarctic Circumpolar Current (ACC) is widely believed to mark the beginning of the evolution of Antarctic fish. The ACC moves in a clockwise northeast direction, and can be up to wide. This current formed 25–22 million years ago, and thermally isolated the Southern Ocean by separating it from the warm subtropical gyres to the north.
During the mid-Tertiary period, a species crash in the Southern Ocean opened up wide range of empty niches to colonize. Despite the hemoglobin-less mutants being less fit, the lack of competition allowed even the mutants to leave descendants that colonized empty habitats and evolved compensations for their mutations. Later, the periodic openings of fjords created habitats that were colonized by a few individuals. These conditions may have also allowed for the loss of myoglobin.
Loss of hemoglobin
The loss of hemoglobin was initially believed to be an adaptation to the extreme cold, as the lack of hemoglobin and red blood cells decreases blood viscosity, which is an adaptation that has been seen in species adapted to cold climates. In refuting this original hypothesis, previous analysis has proposed that the lack of hemoglobin, while not lethal, is not adaptive. Any adaptive advantages incurred by reduced blood viscosity are outweighed by the fact that icefish must pump much more blood per unit of time to make up for the reduced oxygen carrying capacity of their blood. The high blood volume of icefish is itself evidence that the loss of hemoglobin and myoglobin was not advantageous for the ancestor of the icefish. Their unusual cardiovascular physiology, including large heart, high blood volume, increased mitochondrial density, and extensive microvasculature, suggests that icefish have had to evolve ways of coping with the impairment of their oxygen binding and transport systems.
Recent research by Corliss et al. (2019) claims that the loss of hemoglobin has adaptive value. Iron is a limiting nutrient in the environments inhabited by the icefish. By no longer synthesizing hemoglobin, they claim that icefish are minimizing endogenous iron use. To demonstrate this, they obtained retinal samples of Champsocephalus gunnari and stained them to detect hemoglobin alpha 3'f. They found expression of hemoglobin alpha 3'f within the retinal vasculature of Champsocephalus gunnari, demonstrating for the first time that there is limited transcription and translation of a hemoglobin gene fragment within an icefish. Because this fragment of hemoglobin does not contain any iron binding sites, the finding suggests that hemoglobin was selected against to conserve iron.
Loss of myoglobin
Phylogenetic relationships indicate that the nonexpression of myoglobin in cardiac tissue has evolved at least four discrete times. This repeated loss suggests that cardiac myoglobin may be vestigial or even detrimental to icefish. Sidell and O'Brien (2006) investigated this possibility. First, they performed a test using stopped flow spectrometry. They found that across all temperatures, oxygen binds and dissociates faster from icefish than it does from mammalian myoglobin. However, when they repeated the test with each organism at a temperature that accurately reflected its native environment, the myoglobin performance was roughly equivalent between icefish and mammals. So, they concluded that icefish myoglobin is neither more nor less functional than the myoglobin in other clades. This means that myoglobin is unlikely to have been selected against. The same researchers then performed a test in which they selectively inhibited cardiac myoglobin in icefish with natural myoglobin expression. They found that icefish species that naturally lack cardiac myoglobin performed better without myoglobin than did fish that naturally express cardiac myoglobin. This finding suggests that fish without cardiac myoglobin have undergone compensatory adaptation.
Reason for trait fix
The Southern Ocean is an atypical environment. To begin with, the Southern Ocean has been characterized by extremely cold but stable temperatures for the past 10–14 million years. These cold temperatures, which allow for higher water oxygen content, combined with a high degree of vertical mixing in these waters, means oxygen availability in Antarctic waters is unusually high. The loss of hemoglobin and myoglobin would have negative consequences in warmer environments. The stability in temperature is also "lucky", as strong fluctuations in temperature would create a more stressful environment that would likely weed out individuals with deleterious mutations. Although most research suggests that the loss of hemoglobin in icefish was a neutral or maladaptive trait that arose due to a random evolutionary event, some researchers have also suggested that the loss of hemoglobin might be tied to a necessary adaptation for the icefish. Most animals require iron for hemoglobin production, and iron is often limited in ocean environments. Through hemoglobin loss, icefish may minimize their iron requirements. This minimization could have helped the icefish to survive 8.5 million years ago when Arctic diversity plummeted dramatically.
Cardiovascular physiology
The key to solving this conundrum is to consider the other functions that both hemoglobin and myoglobin perform. While emphasis is often placed and understandably so on the importance of hemoglobin and myoglobin in oxygen delivery and use, recent studies have found that both proteins are actually also involved in the process of breaking down nitric oxide. This means that when icefish lost hemoglobin and myoglobin, it did not just mean a decreased ability to transport oxygen, but it also meant that total nitric oxide levels were elevated. Nitric oxide plays a role in regulating various cardiovascular processes in icefish, such as the dilation of branchial vasculature, cardiac stroke volume, and power output. The presence of nitric oxide also can increase angiogenesis, mitochondrial biogenesis, and cause muscle hypertrophy; all of these traits are characteristics of icefish. The similarity between nitric oxide-mediated trait expression and the unusual cardiovascular traits of icefish suggests that while these abnormal traits have evolved over time, much of these traits were simply an immediate physiological response to heightened levels of nitric oxide, which may in turn have led to a process of homeostatic evolution. In addition, the heightened levels of nitric oxide that followed as an inevitable consequence of the loss of hemoglobin and myoglobin may have actually provided an automatic compensation, allowing for the fish to make up for the hit to their oxygen transport system and thereby providing a grace period of the fixation of these less than desirable traits.
| Biology and health sciences | Acanthomorpha | Animals |
693197 | https://en.wikipedia.org/wiki/%CE%9C%20operator | Μ operator | In computability theory, the μ-operator, minimization operator, or unbounded search operator searches for the least natural number with a given property. Adding the μ-operator to the primitive recursive functions makes it possible to define all computable functions.
Definition
Suppose that R(y, x1, ..., xk) is a fixed (k+1)-ary relation on the natural numbers. The μ-operator "μy", in either the unbounded or bounded form, is a "number theoretic function" defined from the natural numbers to the natural numbers. However, "μy" contains a predicate over the natural numbers, which can be thought of as a condition that evaluates to true when the predicate is satisfied and false when it is not.
The bounded μ-operator appears earlier in Kleene (1952) Chapter IX Primitive Recursive Functions, §45 Predicates, prime factor representation as:
"" (p. 225)
Stephen Kleene notes that any of the six inequality restrictions on the range of the variable y is permitted, i.e. y < z, y ≤ z, w < y < z, w < y ≤ z, w ≤ y < z and w ≤ y ≤ z. "When the indicated range contains no y such that R(y) [is "true"], the value of the "μy" expression is the cardinal number of the range" (p. 226); this is why the default "z" appears in the definition above. As shown below, the bounded μ-operator "μyy<z" is defined in terms of two primitive recursive functions called the finite sum Σ and finite product Π, a predicate function that "does the test" and a representing function that converts {t, f} to {0, 1}.
In Chapter XI §57 General Recursive Functions, Kleene defines the unbounded μ-operator over the variable y in the following manner,
"" (p. 279, where "" means "there exists a y such that...")
In this instance R itself, or its representing function, delivers 0 when it is satisfied (i.e. delivers true); the function then delivers the number y. No upper bound exists on y, hence no inequality expressions appear in its definition.
For a given R(y) the unbounded μ-operator μyR(y) (note no requirement for "" ) is a partial function. Kleene makes it as a total function instead (cf. p. 317):
The total version of the unbounded μ-operator is studied in higher-order reverse mathematics () in the following form:
where the superscripts mean that n is zeroth-order, f is first-order, and μ is second-order. This axiom gives rise to the Big Five system ACA0 when combined with the usual base theory of higher-order reverse mathematics.
Properties
(i) In context of the primitive recursive functions, where the search variable y of the μ-operator is bounded, e.g. y < z in the formula below, if the predicate R is primitive recursive (Kleene Proof #E p. 228), then
μyy<zR(y, x1, ..., xn) is a primitive recursive function.
(ii) In the context of the (total) recursive functions, where the search variable y is unbounded but guaranteed to exist for all values xi of the total recursive predicate R's parameters,
(x1),...,(xn) R(y, xi, ..., xn) implies that μyR(y, xi, ..., xn) is a total recursive function.
Here (xi) means "for all xi" and means "there exists at least one value of y such that..." (cf Kleene (1952) p. 279.)
then the five primitive recursive operators plus the unbounded-but-total μ-operator give rise to what Kleene called the "general" recursive functions (i.e. total functions defined by the six recursion operators).
(iii) In the context of the partial recursive functions: Suppose that the relation R holds if and only if a partial recursive function converges to zero. And suppose that that partial recursive function converges (to something, not necessarily zero) whenever μyR(y, x1, ..., xk) is defined and y is μyR(y, x1, ..., xk) or smaller. Then the function μyR(y, x1, ..., xk) is also a partial recursive function.
The μ-operator is used in the characterization of the computable functions as the μ recursive functions.
In constructive mathematics, the unbounded search operator is related to Markov's principle.
Examples
Example 1: The bounded μ-operator is a primitive recursive function
In the following x represents the string xi, ..., xn.
The bounded μ-operator can be expressed rather simply in terms of two primitive recursive functions (hereafter "prf") that also are used to define the CASE function—the product-of-terms Π and the sum-of-terms Σ (cf Kleene #B page 224). (As needed, any boundary for the variable such as s ≤ t or t < z, or 5 < x < 17 etc. is appropriate). For example:
Πs≤t fs(x, s) = f0(x, 0) × f1(x, 1) × ... × ft(x, t)
Σt<z gt(x, t) = g0(x, 0) + g1(x, 1) + ... + gz-1(x, z-1)
Before we proceed we need to introduce a function ψ called "the representing function" of predicate R. Function ψ is defined from inputs (t = "truth", f = "falsity") to outputs (0, 1) (note the order!). In this case the input to ψ. i.e. {t, f}. is coming from the output of R:
ψ(R = t) = 0
ψ(R = f) = 1
Kleene demonstrates that μyy<zR(y) is defined as follows; we see the product function Π is acting like a Boolean OR operator, and the sum Σ is acting somewhat like a Boolean AND but is producing {Σ≠0, Σ=0} rather than just {1, 0}:
μyy<zR(y) = Σt<zΠs≤t ψ(R(x, t, s)) =
[ψ(x, 0, 0)] +
[ψ(x, 1, 0) × ψ(x, 1, 1)] +
[ψ(x, 2, 0) × ψ(x, 2, 1) × ψ(x, 2, 2)] +
... +
[ψ(x, z-1, 0) × ψ(x, z-1, 1) × ψ(x, z-1, 2) × . . . × ψ (x, z-1, z-1)]
Note that Σ is actually a primitive recursion with the base Σ(x, 0) = 0 and the induction step Σ(x, y+1) = Σ(x, y) + Π( x, y). The product Π is also a primitive recursion with base step Π(x, 0) = ψ(x, 0) and induction step Π(x, y+1) = Π(x, y) × ψ(x, y+1).
The equation is easier if observed with an example, as given by Kleene. He just made up the entries for the representing function ψ(R(y)). He designated the representing functions χ(y) rather than ψ(x, y):
Example 2: The unbounded μ-operator is not primitive-recursive
The unbounded μ-operator—the function μy—is the one commonly defined in the texts. But the reader may wonder why the unbounded μ-operator is searching for a function R(x, y) to yield zero, rather than some other natural number.
In a footnote Minsky does allow his operator to terminate when the function inside produces a match to the parameter "k"; this example is also useful because it shows another author's format:
"For μt[φ(t) = k]" (p. 210)
The reason for zero is that the unbounded operator μy will be defined in terms of the function "product" Π with its index y allowed to "grow" as the μ-operator searches. As noted in the example above, the product Πx<y of a string of numbers ψ(x, 0) *, ..., * ψ(x, y) yields zero whenever one of its members ψ(x, i) is zero:
Πs<y = ψ(x, 0) * , ..., * ψ(x, y) = 0
if any ψ(x, i) = 0 where 0≤i≤s. Thus the Π is acting like a Boolean AND.
The function μy produces as "output" a single natural number y = {0, 1, 2, 3, ...}. However, inside the operator one of a couple "situations" can appear: (a) a "number-theoretic function" χ that produces a single natural number, or (b) a "predicate" R that produces either {t = true, f = false}. (And, in the context of partial recursive functions Kleene later admits a third outcome: "μ = undecided".)
Kleene splits his definition of the unbounded μ-operator to handle the two situations (a) and (b). For situation (b), before the predicate R(x, y) can serve in an arithmetic capacity in the product Π, its output {t, f} must first be "operated on" by its representing function χ to yield {0, 1}. And for situation (a) if one definition is to be used then the number theoretic function χ must produce zero to "satisfy" the μ-operator. With this matter settled, he demonstrates with single "Proof III" that either types (a) or (b) together with the five primitive recursive operators yield the (total) recursive functions, with this proviso for a total function:
For all parameters x, a demonstration must be provided to show that a y exists that satisfies (a) μyψ(x, y) or (b) μyR(x, y).
Kleene also admits a third situation (c) that does not require the demonstration of "for all x a y exists such that ψ(x, y)." He uses this in his proof that more total recursive functions exist than can be enumerated; c.f. footnote Total function demonstration.
Kleene's proof is informal and uses an example similar to the first example, but first he casts the μ-operator into a different form that uses the "product-of-terms" Π operating on function χ that yields a natural number n, which can be any natural number, and 0 in the instance when the u-operator's test is "satisfied".
The definition recast with the Π-function:
μyy<zχ(y) =
(i): π(x, y) = Πs<yχ(x, s)
(ii): φ(x) = τ(π(x, y), π(x, y' ), y)
(iii): τ(z' , 0, y) = y ;τ(u, v, w) is undefined for u = 0 or v > 0.
This is subtle. At first glance the equations seem to be using primitive recursion. But Kleene has not provided us with a base step and an induction step of the general form:
base step: φ(0, x) = φ(x)
induction step: φ(0, x) = ψ(y, φ(0,x), x)
To see what is going on, we first have to remind ourselves that we have assigned a parameter (a natural number) to every variable xi. Second, we do see a successor-operator at work iterating y (i.e. the y' ). And third, we see that the function μy y<zχ(y, x) is just producing instances of χ(y,x) i.e. χ(0,x), χ(1,x), ... until an instance yields 0. Fourth, when an instance χ(n, x) yields 0 it causes the middle term of τ, i.e. v = π(x, y' ) to yield 0. Finally, when the middle term v = 0, μyy<zχ(y) executes line (iii) and "exits". Kleene's presentation of equations (ii) and (iii) have been exchanged to make this point that line (iii) represents an exit—an exit taken only when the search successfully finds a y to satisfy χ(y) and the middle product-term π(x, y' ) is 0; the operator then terminates its search with τ(z' , 0, y) = y.
τ(π(x, y), π(x, y' ), y), i.e.:
τ(π(x, 0), π(x, 1), 0),
τ(π(x, 1), π(x, 2), 1)
τ(π(x, 2), π(x, 3), 2)
τ(π(x, 3), π(x, 4), 3)
... until a match occurs at y=n and then:
τ(z' , 0, y) = τ(z' , 0, n) = n and the μ-operator's search is done.
For the example Kleene "...consider[s] any fixed values of (xi, ..., xn) and write[s] simply 'χ(y)' for 'χ(xi, ..., xn), y)'":
Example 3: Definition of the unbounded μ-operator in terms of an abstract machine
Both Minsky (1967) p. 21 and Boolos-Burgess-Jeffrey (2002) p. 60-61 provide definitions of the μ-operator as an abstract machine; see footnote Alternative definitions of μ.
The following demonstration follows Minsky without the "peculiarity" mentioned in the footnote. The demonstration will use a "successor" counter machine model closely related to the Peano Axioms and the primitive recursive functions. The model consists of (i) a finite state machine with a TABLE of instructions and a so-called 'state register' that we will rename "the Instruction Register" (IR), (ii) a few "registers" each of which can contain only a single natural number, and (iii) an instruction set of four "commands" described in the following table:
In the following, the symbolism " [ r ] " means "the contents of", and " →r " indicates an action with respect to register r.
The algorithm for the minimization operator μy[φ(x, y)] will, in essence, create a sequence of instances of the function φ(x, y) as the value of parameter y (a natural number) increases; the process will continue (see Note † below) until a match occurs between the output of function φ(x, y) and some pre-established number (usually 0). Thus the evaluation of φ(x, y) requires, at the outset, assignment of a natural number to each of its variables x and an assignment of a "match-number" (usually 0) to a register "w", and a number (usually 0) to register y.
Note †: The unbounded μ-operator will continue this attempt-to-match process ad infinitum or until a match occurs. Thus the "y" register must be unbounded -- it must be able to "hold" a number of arbitrary size. Unlike a "real" computer model, abstract machine models allow this. In the case of a bounded μ-operator, a lower-bounded μ-operator would start with the contents of y set to a number other than zero. An upper-bounded μ-operator would require an additional register "ub" to contain the number that represents the upper bound plus an additional comparison operation; an algorithm could provide for both lower- and upper bounds.
In the following we are assuming that the Instruction Register (IR) encounters the μy "routine" at instruction number "n". Its first action will be to establish a number in a dedicated "w" register—an "example of" the number that function φ(x, y) must produce before the algorithm can terminate (classically this is the number zero, but see the footnote about the use of numbers other than zero). The algorithm's next action at instructiton "n+1" will be to clear the "y" register -- "y" will act as an "up-counter" that starts from 0. Then at instruction "n+2" the algorithm evaluates its function φ(x, y) -- we assume this takes j instructions to accomplish—and at the end of its evaluation φ(x, y) deposits its output in register "φ". At the (n+j+3)rd instruction the algorithm compares the number in the "w" register (e.g. 0) to the number in the "φ" register—if they are the same the algorithm has succeeded and it escapes through exit; otherwise it increments the contents of the "y" register and loops back with this new y-value to test function φ(x, y) again.
| Mathematics | Computability theory | null |
693334 | https://en.wikipedia.org/wiki/Foundation%20%28engineering%29 | Foundation (engineering) | In engineering, a foundation is the element of a structure which connects it to the ground or more rarely, water (as with floating structures), transferring loads from the structure to the ground. Foundations are generally considered either shallow or deep. Foundation engineering is the application of soil mechanics and rock mechanics (geotechnical engineering) in the design of foundation elements of structures.
Purpose
Foundations provide the structure's stability from the ground:
To distribute the weight of the structure over a large area in order to avoid overloading the underlying soil (possibly causing unequal settlement).
To anchor the structure against natural forces including earthquakes, floods, droughts, frost heaves, tornadoes and wind.
To provide a level surface for construction.
To anchor the structure deeply into the ground, increasing its stability and preventing overloading.
To prevent lateral movements of the supported structure (in some cases).
Requirements of a good foundation
The design and the construction of a well-performing foundation must possess some basic requirements:
The design and the construction of the foundation is done such that it can sustain as well as transmit the dead and the imposed loads to the soil. This transfer has to be carried out without resulting in any form of settlement that can cause stability issues for the structure.
Differential settlements can be avoided by having a rigid base for the foundation. These issues are more pronounced in areas where the superimposed loads are not uniform in nature.
Based on the soil and area it is recommended to have a deeper foundation so that it can guard any form of damage or distress. These are mainly caused due to the problem of shrinkage and swelling because of temperature changes.
The location of the foundation chosen must be an area that is not affected or influenced by future works or factors.
Historic types
Earthfast or post in ground construction
Buildings and structures have a long history of being built with wood in contact with the ground. Post in ground construction may technically have no foundation. Timber pilings were used on soft or wet ground even below stone or masonry walls. In marine construction and bridge building a crisscross of timbers or steel beams in concrete is called grillage.
Padstones
Perhaps the simplest foundation is the padstone, a single stone which both spreads the weight on the ground and raises the timber off the ground. Staddle stones are a specific type of padstone.
Stone foundations
Dry stone and stones laid in mortar to build foundations are common in many parts of the world. Dry laid stone foundations may have been painted with mortar after construction. Sometimes the top, visible course of stone is hewn, quarried stones. Besides using mortar, stones can also be put in a gabion. One disadvantage is that if using regular steel rebars, the gabion would last much less long than when using mortar (due to rusting). Using weathering steel rebars could reduce this disadvantage somewhat.
Rubble-trench foundations
Rubble trench foundations are a shallow trench filled with rubble or stones. These foundations extend below the frost line and may have a drain pipe which helps groundwater drain away. They are suitable for soils with a capacity of more than 10 tonnes/m2 (2,000 pounds per square foot).
Gallery of shallow foundation types
Modern types
Shallow foundations
Often called footings, are usually embedded about a meter or so into soil. One common type is the spread footing which consists of strips or pads of concrete (or other materials) which extend below the frost line and transfer the weight from walls and columns to the soil or bedrock.
Another common type of shallow foundation is the slab-on-grade foundation where the weight of the structure is transferred to the soil through a concrete slab placed at the surface. Slab-on-grade foundations can be reinforced mat slabs, which range from 25 cm to several meters thick, depending on the size of the building, or post-tensioned slabs, which are typically at least 20 cm for houses, and thicker for heavier structures.
Another way to install ready-to-build foundations that is more environmentally friendly is to use screw piles. Screw pile installations have also extended to residential applications, with many homeowners choosing a screw pile foundation over other options. Some common applications for helical pile foundations include wooden decks, fences, garden houses, pergolas, and carports.
Deep foundations
Used to transfer the load of a structure down through the upper weak layer of topsoil to the stronger layer of subsoil below. There are different types of deep footings including impact driven piles, drilled shafts, caissons, screw piles, geo-piers and earth-stabilized columns. The naming conventions for different types of footings vary between different engineers. Historically, piles were wood, later steel, reinforced concrete, and pre-tensioned concrete.
Monopile foundation
A type of deep foundation which uses a single, generally large-diameter, structural element embedded into the earth to support all the loads (weight, wind, etc.) of a large above-surface structure.
Many monopile foundations have been used in recent years for economically constructing fixed-bottom offshore wind farms in shallow-water subsea locations.
For example, a single wind farm off the coast of England went online in 2008 with over 100 turbines, each mounted on a 4.74-meter-diameter monopile footing in ocean depths up to 16 meters of water.
Floating\barge
A floating foundation is one that sits on a body of water, rather than dry land. This type of foundation is used for some bridges and floating buildings.
Design
Foundations are designed to have an adequate load capacity depending on the type of subsoil/rock supporting the foundation by a geotechnical engineer, and the footing itself may be designed structurally by a structural engineer. The primary design concerns are settlement and bearing capacity. When considering settlement, total settlement and differential settlement is normally considered. Differential settlement is when one part of a foundation settles more than another part. This can cause problems to the structure which the foundation is supporting. Expansive clay soils can also cause problems.
| Technology | Architectural elements | null |
693623 | https://en.wikipedia.org/wiki/Galactic%20plane | Galactic plane | The galactic plane is the plane on which the majority of a disk-shaped galaxy's mass lies. The directions perpendicular to the galactic plane point to the galactic poles. In actual usage, the terms galactic plane and galactic poles usually refer specifically to the plane and poles of the Milky Way, in which Planet Earth is located.
Some galaxies are irregular and do not have any well-defined disk. Even in the case of a barred spiral galaxy like the Milky Way, defining the galactic plane is slightly imprecise and arbitrary since the stars are not perfectly coplanar. In 1959, the IAU defined the position of the Milky Way's north galactic pole as exactly RA = , Dec = in the then-used B1950 epoch; in the currently-used J2000 epoch, after precession is taken into account, its position is RA , Dec . This position is in Coma Berenices, near the bright star Arcturus; likewise, the south galactic pole lies in the constellation Sculptor.
The zero of longitude of galactic coordinates was also defined in 1959 to be at position angle 123° from the north celestial pole. Thus the zero longitude point on the galactic equator was at , (B1950) or , (J2000), and its J2000 position angle is 122.932°. The Galactic Center is located at position angle 31.72° (B1950) or 31.40° (J2000) east of north.
| Physical sciences | Basics_2 | Astronomy |
2404348 | https://en.wikipedia.org/wiki/Magnetization | Magnetization | In classical electromagnetism, magnetization is the vector field that expresses the density of permanent or induced magnetic dipole moments in a magnetic material. Accordingly, physicists and engineers usually define magnetization as the quantity of magnetic moment per unit volume.
It is represented by a pseudovector M. Magnetization can be compared to electric polarization, which is the measure of the corresponding response of a material to an electric field in electrostatics.
Magnetization also describes how a material responds to an applied magnetic field as well as the way the material changes the magnetic field, and can be used to calculate the forces that result from those interactions.
The origin of the magnetic moments responsible for magnetization can be either microscopic electric currents resulting from the motion of electrons in atoms, or the spin of the electrons or the nuclei. Net magnetization results from the response of a material to an external magnetic field.
Paramagnetic materials have a weak induced magnetization in a magnetic field, which disappears when the magnetic field is removed. Ferromagnetic and ferrimagnetic materials have strong magnetization in a magnetic field, and can be magnetized to have magnetization in the absence of an external field, becoming a permanent magnet. Magnetization is not necessarily uniform within a material, but may vary between different points.
Definition
The magnetization field or M-field can be defined according to the following equation:
Where is the elementary magnetic moment and is the volume element; in other words, the M-field is the distribution of magnetic moments in the region or manifold concerned. This is better illustrated through the following relation:
where m is an ordinary magnetic moment and the triple integral denotes integration over a volume. This makes the M-field completely analogous to the electric polarisation field, or P-field, used to determine the electric dipole moment p generated by a similar region or manifold with such a polarization:
where is the elementary electric dipole moment.
Those definitions of P and M as a "moments per unit volume" are widely adopted, though in some cases they can lead to ambiguities and paradoxes.
The M-field is measured in amperes per meter (A/m) in SI units.
In Maxwell's equations
The behavior of magnetic fields (B, H), electric fields (E, D), charge density (ρ), and current density (J) is described by Maxwell's equations. The role of the magnetization is described below.
Relations between B, H, and M
The magnetization defines the auxiliary magnetic field H as
(SI)
(Gaussian system)
which is convenient for various calculations. The vacuum permeability μ0 is, approximately, ).
A relation between M and H exists in many materials. In diamagnets and paramagnets, the relation is usually linear:
where χ is called the volume magnetic susceptibility, and μ is called the magnetic permeability of the material. The magnetic potential energy per unit volume (i.e. magnetic energy density) of the paramagnet (or diamagnet) in the magnetic field is:
the negative gradient of which is the magnetic force on the paramagnet (or diamagnet) per unit volume (i.e. force density).
In diamagnets () and paramagnets (), usually , and therefore .
In ferromagnets there is no one-to-one correspondence between M and H because of magnetic hysteresis.
Magnetic polarization
Alternatively to the magnetization, one can define the magnetic polarization, (often the symbol is used, not to be confused with current density).
(SI).
This is by direct analogy to the electric polarization, .
The magnetic polarization thus differs from the magnetization by a factor of :
(SI).
Whereas magnetization is given with the unit ampere/meter, the magnetic polarization is given with the unit tesla.
Magnetization current
The magnetization M makes a contribution to the current density J, known as the magnetization current.
and for the bound surface current:
so that the total current density that enters Maxwell's equations is given by
where Jf is the electric current density of free charges (also called the free current), the second term is the contribution from the magnetization, and the last term is related to the electric polarization P.
Magnetostatics
In the absence of free electric currents and time-dependent effects, Maxwell's equations describing the magnetic quantities reduce to
These equations can be solved in analogy with electrostatic problems where
In this sense −∇⋅M plays the role of a fictitious "magnetic charge density" analogous to the electric charge density ρ; (see also demagnetizing field).
Dynamics
The time-dependent behavior of magnetization becomes important when considering nanoscale and nanosecond timescale magnetization. Rather than simply aligning with an applied field, the individual magnetic moments in a material begin to precess around the applied field and come into alignment through relaxation as energy is transferred into the lattice.
Reversal
Magnetization reversal, also known as switching, refers to the process that leads to a 180° (arc) re-orientation of the magnetization vector with respect to its initial direction, from one stable orientation to the opposite one. Technologically, this is one of the most important processes in magnetism that is linked to the magnetic data storage process such as used in modern hard disk drives. As it is known today, there are only a few possible ways to reverse the magnetization of a metallic magnet:
an applied magnetic field
spin injection via a beam of particles with spin
magnetization reversal by circularly polarized light; i.e., incident electromagnetic radiation that is circularly polarized
Demagnetization
Demagnetization is the reduction or elimination of magnetization. One way to do this is to heat the object above its Curie temperature, where thermal fluctuations have enough energy to overcome exchange interactions, the source of ferromagnetic order, and destroy that order. Another way is to pull it out of an electric coil with alternating current running through it, giving rise to fields that oppose the magnetization.
One application of demagnetization is to eliminate unwanted magnetic fields. For example, magnetic fields can interfere with electronic devices such as cell phones or computers, and with machining by making cuttings cling to their parent.
| Physical sciences | Magnetostatics | Physics |
2408688 | https://en.wikipedia.org/wiki/Nucleophilic%20aromatic%20substitution | Nucleophilic aromatic substitution | A nucleophilic aromatic substitution (SNAr) is a substitution reaction in organic chemistry in which the nucleophile displaces a good leaving group, such as a halide, on an aromatic ring. Aromatic rings are usually nucleophilic, but some aromatic compounds do undergo nucleophilic substitution. Just as normally nucleophilic alkenes can be made to undergo conjugate substitution if they carry electron-withdrawing substituents, so normally nucleophilic aromatic rings also become electrophilic if they have the right substituents.This reaction differs from a common SN2 reaction, because it happens at a trigonal carbon atom (sp2 hybridization). The mechanism of SN2 reaction does not occur due to steric hindrance of the benzene ring. In order to attack the C atom, the nucleophile must approach in line with the C-LG (leaving group) bond from the back, where the benzene ring lies. It follows the general rule for which SN2 reactions occur only at a tetrahedral carbon atom.
The SN1 mechanism is possible but very unfavourable unless the leaving group is an exceptionally good one. It would involve the unaided loss of the leaving group and the formation of an aryl cation. In the SN1 reactions all the cations employed as intermediates were planar with an empty p orbital. This cation is planar but the p orbital is full (it is part of the aromatic ring) and the empty orbital is an sp2 orbital outside the ring.
Nucleophilic aromatic substitution mechanisms
Aromatic rings undergo nucleophilic substitution by several pathways.
SNAr (addition-elimination) mechanism
aromatic SN1 mechanism encountered with diazonium salts
benzyne mechanism (E1cB-AdN)
free radical SRN1 mechanism
ANRORC mechanism
Vicarious nucleophilic substitution
The SNAr mechanism is the most important of these. Electron withdrawing groups activate the ring towards nucleophilic attack. For example if there are nitro functional groups positioned ortho or para to the halide leaving group, the SNAr mechanism is favored.
SNAr reaction mechanism
The following is the reaction mechanism of a nucleophilic aromatic substitution of 2,4-dinitrochlorobenzene (1) in a basic solution in water.
Since the nitro group is an activator toward nucleophilic substitution, and a meta director, it is able to stabilize the additional electron density (via resonance) when the aromatic compound is attacked by the hydroxide nucleophile. The resulting intermediate, named the Meisenheimer complex (2a), the ipso carbon is temporarily bonded to the hydroxyl group. This Meisenheimer complex is extra stabilized by the additional electron-withdrawing nitro group (2b).
In order to return to a lower energy state, either the hydroxyl group leaves, or the chloride leaves. In solution, both processes happen. A small percentage of the intermediate loses the chloride to become the product (2,4-dinitrophenol, 3), while the rest return to the reactant (1). Since 2,4-dinitrophenol is in a lower energy state, it will not return to form the reactant, so after some time has passed, the reaction reaches chemical equilibrium that favors the 2,4-dinitrophenol, which is then deprotonated by the basic solution (4).
The formation of the resonance-stabilized Meisenheimer complex is slow because the loss of aromaticity due to nucleophilic attack results in a higher-energy state. By the same coin, the loss of the chloride or hydroxide is fast, because the ring regains aromaticity. Recent work indicates that, sometimes, the Meisenheimer complex is not always a true intermediate but may be the transition state of a 'frontside SN2' process, particularly if stabilization by electron-withdrawing groups is not very strong. A 2019 review argues that such 'concerted SNAr' reactions are more prevalent than previously assumed.
Aryl halides cannot undergo the classic 'backside' SN2 reaction. The carbon-halogen bond is in the plane of the ring because the carbon atom has a trigonal planar geometry. Backside attack is blocked and this reaction is therefore not possible. An SN1 reaction is possible but very unfavourable. It would involve the unaided loss of the leaving group and the formation of an aryl cation. The nitro group is the most commonly encountered activating group, other groups are the cyano and the acyl group. The leaving group can be a halogen or a sulfide. With increasing electronegativity the reaction rate for nucleophilic attack increases. This is because the rate-determining step for an SNAr reaction is attack of the nucleophile and the subsequent breaking of the aromatic system; the faster process is the favourable reforming of the aromatic system after loss of the leaving group. As such, the following pattern is seen with regard to halogen leaving group ability for SNAr: F > Cl ≈ Br > I (i.e. an inverted order to that expected for an SN2 reaction). If looked at from the point of view of an SN2 reaction this would seem counterintuitive, since the C-F bond is amongst the strongest in organic chemistry, when indeed the fluoride is the ideal leaving group for an SNAr due to the extreme polarity of the C-F bond. Nucleophiles can be amines, alkoxides, sulfides and stabilized carbanions.
Nucleophilic aromatic substitution reactions
Some typical substitution reactions on arenes are listed below.
In the Bamberger rearrangement N-phenylhydroxylamines rearrange to 4-aminophenols. The nucleophile is water.
The Smiles rearrangement is the intramolecular version of this reaction type.
Nucleophilic aromatic substitution is not limited to arenes, however; the reaction takes place even more readily with heteroarenes. Pyridines are especially reactive when substituted in the aromatic ortho position or aromatic para position because then the negative charge is effectively delocalized at the nitrogen position. One classic reaction is the Chichibabin reaction (Aleksei Chichibabin, 1914) in which pyridine is reacted with an alkali-metal amide such as sodium amide to form 2-aminopyridine.
In the compound methyl 3-nitropyridine-4-carboxylate, the meta nitro group is actually displaced by fluorine with cesium fluoride in DMSO at 120 °C.
Although the Sandmeyer reaction of diazonium salts and halides is formally a nucleophilic substitution, the reaction mechanism is in fact radical.
Asymmetric nucleophilic aromatic substitution
With carbon nucleophiles such as 1,3-dicarbonyl compounds the reaction has been demonstrated as a method for the asymmetric synthesis of chiral molecules. First reported in 2005, the organocatalyst (in a dual role with that of a phase transfer catalyst) is derived from cinchonidine (benzylated at N and O).
| Physical sciences | Organic reactions | Chemistry |
2409351 | https://en.wikipedia.org/wiki/Armor-piercing%20bullet | Armor-piercing bullet | Armor-piercing bullets for rifle and handgun cartridges are designed to penetrate ballistic armor and protective shields intended to stop or deflect conventional bullets. Although bullet design is an important factor with regard to armor penetration, the ability of any given projectile to penetrate ballistic armor increases with increasing velocity. Rifle cartridges typically discharge bullets at higher muzzle velocity than handgun cartridges due to larger propellant charge. However, even the same cartridge (one that is interchangeable between specific rifles and handguns) fired from a rifle will, in almost all common cases, have a higher velocity than when fired from a handgun. This is due to the longer period of acceleration available within the longer gun barrel of rifles, which allow adequate time for the propellant to fully ignite before the projectile exits the barrel. For this reason, bullets fired from rifles may be more capable of piercing armor than similar or identical bullets fired from handguns. In addition, a small-caliber bullet has higher sectional density than a larger-caliber bullet of the same weight, and thus is more capable of defeating body armor.
Design
Rifle bullets
Armor-piercing bullets typically contain a hardened steel, tungsten, or tungsten carbide penetrator encased within a copper or cupronickel jacket, similar to the jacket which would surround lead in a conventional projectile. The penetrator is a pointed mass of high-density material designed to retain its shape and carry the maximum possible amount of energy as deeply as possible into the target. The entire projectile is not normally made of the same material as the penetrator because the hard metals of good penetrators would damage the barrel of the gun firing the bullet. Impact velocity of the copper jacket may temporarily soften the face of the armor and cushion the impact to avoid breaking the brittle penetrator. The penetrator then slides out of the jacket to continue forward through the armor.
Examples
Handgun bullets
Handgun bullets made entirely of lead have less penetration ability than jacketed bullets at similar velocity. In the 1930s, Western Cartridge Company introduced .38 Special ammunition capable of firing a copper-tipped lead-alloy bullet at per second to penetrate sheet-metal automobile doors. As higher velocity handgun cartridges became available and jacketed bullets became more common in handgun cartridges, armor penetration was improved with thicker bullet jackets or bullets made entirely of jacket material like copper or brass. Later designs used penetrator cores similar to rifle designs. Some of these bullets were coated with Teflon to reduce their tendency to ricochet off glass or sheet metal.
United States
In 1986 United States law initially defined armor-piercing bullets to exempt rifle ammunition:
A projectile or projectile core which may be used in a handgun and which is constructed entirely (excluding the presence of traces of other substances) from one or a combination of tungsten alloys, steel, iron, brass, bronze, beryllium copper, or depleted uranium; or
A full jacketed projectile larger than .22 caliber designed and intended for use in a handgun and whose jacket has a weight of more than 25 percent of the total weight of the projectile.
Subsequent regulations requiring green bullets encouraged replacing lead core bullets with M855A1 military bullets with a copper jacket over a steel core, or hunting bullets of solid copper or brass.
Sources
Bullets | Technology | Ammunition | null |
15124887 | https://en.wikipedia.org/wiki/Diphosphane | Diphosphane | Diphosphane, or diphosphine, is an inorganic compound with the chemical formula P2H4. This colourless liquid is one of several binary phosphorus hydrides. It is the impurity that typically causes samples of phosphine to ignite in air.
Properties, preparation, reactions
Diphosphane adopts the gauche conformation (like hydrazine, less symmetrical than shown in the image) with a P−P distance of 2.219 angstroms. It is nonbasic, unstable at room temperature, and spontaneously flammable in air. It is only poorly soluble in water but dissolves in organic solvents. Its 1H NMR spectrum consists of 32 lines resulting from an A2XX'A'2 splitting system.
Diphosphane is produced by the hydrolysis of calcium monophosphide, which can be described as the Ca2+ derivative of . According to an optimized procedure, hydrolysis of 400 g of CaP at −30 °C gives about 20 g of product, slightly contaminated with phosphine.
Reaction of diphosphane with butyllithium affords a variety of condensed polyphosphine compounds.
Organic diphosphanes
A variety of organic derivatives of diphosphane are known, but asymmetric diphosphanes are only stable at cryogenic temperatures. Otherwise, the substituents facily redistribute on the phosphorus centers to give a mixture of products. On the other hand, there appears to be a substantial barrier to chiral inversion.
The central bond is weak, and easily adds substituents.
The simplest synthesis method heats a phosphorus halide and a phosphane:
Ph2PCl + HPPh2 → Ph2P-PPh2 + HCl↑
Alkali metals can replace the hydrogen in that reaction (i.e., a dialkylphosphide), and in some rare cases a dialkylamine can replace the halide. Symmetric diphosphanes are easily prepared by reductive coupling, e.g. tetraphenyldiphosphine from chlorodiphenylphosphine:
2 ClPPh2 + 2 Na → Ph2P−PPh2 + 2 NaCl
Ultraviolet radiation decomposes mercury(II) dialkylphosphides to the metal and a dialkylphosphane.
The methyl compound P2Me4 is prepared by the reduction of Me2P(S)−P(S)Me2, which is produced by methylation of thiophosphoryl chloride with methylmagnesium bromide.
| Physical sciences | Hydrogen compounds | Chemistry |
19074264 | https://en.wikipedia.org/wiki/Amalgam%20%28chemistry%29 | Amalgam (chemistry) | An amalgam is an alloy of mercury with another metal. It may be a liquid, a soft paste or a solid, depending upon the proportion of mercury. These alloys are formed through metallic bonding, with the electrostatic attractive force of the conduction electrons working to bind all the positively charged metal ions together into a crystal lattice structure. Almost all metals can form amalgams with mercury, the notable exceptions being iron, platinum, tungsten, and tantalum. Silver-mercury amalgams are important in dentistry, and gold-mercury amalgam is used in the extraction of gold from ore. Dentistry has used alloys of mercury with metals such as silver, copper, indium, tin and zinc.
Important amalgams
Zinc amalgam
Zinc amalgam finds use in organic synthesis (e.g., for the Clemmensen reduction).
It is the reducing agent in the Jones reductor, used in analytical chemistry. Formerly the zinc plates of dry batteries were amalgamated with a small amount of mercury to prevent deterioration in storage. It is a binary solution (liquid-solid) of mercury and zinc.
Potassium amalgam
For the alkali metals, amalgamation is exothermic, and distinct chemical forms can be identified, such as KHg and KHg2. KHg is a gold-coloured compound with a melting point of 178 °C, and KHg2 a silver-coloured compound with a melting point of 278 °C. These amalgams are very sensitive to air and water, but can be worked with under dry nitrogen. The Hg-Hg distance is around 300 picometres, Hg-K around 358 pm.
Phases K5Hg7 and KHg11 are also known; rubidium, strontium and barium undecamercurides are known and isostructural. Sodium amalgam (NaHg2) has a different structure, with the mercury atoms forming hexagonal layers, and the sodium atoms a linear chain which fits into the holes in the hexagonal layers, but the potassium atom is too large for this structure to work in KHg2.
Sodium amalgam
Sodium amalgam is produced as a byproduct of the chloralkali process and used as an important reducing agent in organic and inorganic chemistry. With water, it decomposes into concentrated sodium hydroxide solution, hydrogen and mercury, which can then return to the chloralkali process anew. If absolutely water-free alcohol is used instead of water, an alkoxide of sodium is produced instead of the alkali solution.
Aluminium amalgam
Aluminium can form an amalgam through a reaction with mercury. Aluminium amalgam may be prepared by either grinding aluminium pellets or wire in mercury, or by allowing aluminium wire or foil to react with a solution of mercuric chloride. This amalgam is used as a reagent to reduce compounds, such as the reduction of imines to amines. The aluminium is the ultimate electron donor, and the mercury serves to mediate the electron transfer.
The reaction itself and the waste from it contain mercury, so special safety precautions and disposal methods are needed. As an environmentally friendlier alternative, hydrides or other reducing agents can often be used to accomplish the same synthetic result. Another environmentally friendly alternative is an alloy of aluminium and gallium which similarly renders the aluminium more reactive by preventing it from forming an oxide layer.
Tin amalgam
Tin amalgam was used in the middle of the 19th century as a reflective mirror coating.
Other amalgams
A variety of amalgams are known that are of interest mainly in the research context.
Ammonium amalgam is a grey, soft, spongy mass discovered in 1808 by Humphry Davy and Jöns Jakob Berzelius. It decomposes readily at room temperature or in contact with water or alcohol:
Thallium amalgam has a freezing point of −58 °C, which is lower than that of pure mercury (−38.8 °C) so it has found a use in low temperature thermometers.
Gold amalgam: Refined gold, when finely ground and brought into contact with mercury where the surfaces of both metals are clean, amalgamates readily and quickly forms alloys ranging from AuHg2 to Au8Hg.
Lead forms an amalgam when filings are mixed with mercury and is also listed as a naturally occurring alloy called leadamalgam in the Nickel–Strunz classification.
Dental amalgam
Dentistry has used alloys of mercury with metals such as silver, copper, indium, tin and zinc. Amalgam is an "excellent and versatile restorative material" and is used in dentistry because it is inexpensive and relatively easy to use and manipulate during placement. It remains soft for a short time so it can be packed to fill any irregular volume, and then forms a hard compound. Amalgam possesses greater longevity when compared to other direct restorative materials, such as composite. However, this difference has decreased with continual development of composite resins.
Amalgam is typically compared to resin-based composites because many applications are similar and many physical properties and costs are comparable.
Dental amalgam has been studied and is generally considered to be safe for humans, though the validity of some studies and their conclusions have been questioned.
In July 2018 the EU, in consideration of the persistent pollution and environmental toxicity of amalgam's mercury, prohibited amalgam for dental treatment of children under 15 years and of pregnant or breastfeeding women.
Use in mining
Mercury has been used in gold and silver mining because of the convenience and the ease with which mercury and the precious metals will amalgamate. In gold placer mining, in which minute specks of gold are washed from sand or gravel deposits, mercury was often used to separate the gold from other heavy minerals.
After all of the practical metal had been taken out from the ore, the mercury was dispensed down a long copper trough, which formed a thin coating of mercury on the exterior. The waste ore was then transferred down the trough, and gold in the waste amalgamated with the mercury. This coating would then be scraped off and refined by evaporation to get rid of the mercury, leaving behind somewhat high-purity gold.
Mercury amalgamation was first used on silver ores with the development of the patio process in Mexico in 1557. There were also additional amalgamation processes that were created for processing silver ores, including pan amalgamation and the Washoe process.
Gold amalgam
Gold extraction (mining)
Gold amalgam has proved effective where gold fines ("flour gold") would not be extractable from ore using hydro-mechanical methods. Large amounts of mercury were used in placer mining, where deposits composed largely of decomposed granite slurry were separated in long runs of "riffle boxes", with mercury dumped in at the head of the run. The amalgam formed is a heavy dull gray solid mass. The use of mercury in 19th century placer mining in California, now prohibited, has caused extensive pollution problems in riverine and estuarine environments, ongoing to this day. Sometimes substantial slugs of amalgam are found in downstream river and creek bottoms by amateur wet-suited miners seeking gold nuggets with the aid of an engine-powered water vacuum/dredge mounted on a float.
Gold extraction (ore processing)
Where stamp mills were used to crush gold-bearing ore to fines, a part of the extraction process involved the use of mercury-wetted copper plates, over which the crushed fines were washed. A periodic scraping and re-mercurizing of the plate resulted in amalgam for further processing.
Gold extraction (retorting)
Amalgam obtained by either process was then heated in a distillation retort, recovering the mercury for reuse and leaving behind the gold. As this released mercury vapors to the atmosphere, the process could induce adverse health effects and long term pollution.
Today, mercury amalgamation has been replaced by other methods to recuperate gold and silver from ore in developed nations. Hazards of mercurial toxic waste have played a major role in the phasing out of the mercury amalgamation processes. Mercury amalgamation is still regularly used by small-scale gold placer miners (often illegally), particularly in developing countries.
Amalgam probe
Mercury salts are, compared to mercury metal and amalgams, highly toxic due to their solubility in water. The presence of these salts in water can be detected with a probe that uses the readiness of mercury ions to form an amalgam with copper. A nitric acid solution of salts under investigation is applied to a piece of copper foil, and any mercury ions present will leave spots of silvery-coloured amalgam. Silver ions leave similar spots but are easily washed away, making this a means of distinguishing silver from mercury.
The redox reaction involved where mercury oxidizes the copper is:
Hg2+ + Cu → Hg + Cu2+.
| Physical sciences | Chemical mixtures: General | null |
10027284 | https://en.wikipedia.org/wiki/Hook%20%28hand%20tool%29 | Hook (hand tool) | A hook is a hand tool used for securing and moving loads. It consists of a round wooden handle with a strong metal hook about long projecting at a right angle from the center of the handle. The appliance is held in a closed fist with the hook projecting between two fingers.
This type of hook is used in many different industries, and has many different names. It may be called a box hook, cargo hook, loading hook, docker's hook when used by longshoremen, and a baling hook, bale hook, or hay hook in the agricultural industry. Other variants exist, such as in forestry, for moving logs, and a type with a long shaft, used by city workers to remove manhole covers.
Smaller hooks may also be used in food processing and transport.
Dockwork
The longshoreman's hook was historically used by longshoremen (stevedores). Before the age of containerization, freight was moved on and off ships with extensive manual labor, and the longshoreman's hook was the basic tool of the dockworker. The hook became an emblem of the longshoreman's profession in the same way that a hammer and anvil are associated with blacksmiths, or the pipe wrench with pipefitters, sprinklerfitters and plumbers. When longshoremen went on strike or retired, it was known as "hanging up the hook" or "slinging the hook", and the newsletter for retired members of the International Longshore and Warehouse Union's Seattle Local is called The Rusty Hook. A longshoreman's hook was often carried by hooking it through the belt.
Longshoremen carried various types of hooks depending on the cargo they would handle. Cargo could come in the form of bales, sacks, barrels, wood crates, or it could be stowed individually in the cargo hold of the ship. The primary function of the hook was to protect the hands of the longshoreman from being injured while handling the cargo. Hooks also improved the reach of the worker and allowed greater strength and handling of the cargo.
Some cargo items are liable to be damaged if pulled at with a longshoreman's hook: hence the "Use No Hooks" warning sign.
A longshoreman's hook looks somewhat intimidating, and as it was also associated with strong, tough dockworkers, it became a commonly used weapon in crime fiction, similar to the ice pick. For example, in an episode of Alfred Hitchcock Presents entitled Shopping for Death, a character is murdered (off-screen) using a longshoreman's hook. It was sometimes used as a weapon and means of intimidation in real life as well; the book Joey the Hit Man: The Autobiography of a Mafia Killer states "One guy who used to work on the docks was called Charlie the Hook. If he didn't like you he would pick you up with his hook." In the 1957 New York drama film Edge of the City, two longshoremen settle their dispute in a deadly baling hook fight. They are also the primary weapon of Spider Splicers in the BioShock series, so named due to their use of the hooks to crawl on ceilings and attack unexpectedly.
Haying
A hay hook is slightly different in design from a longshoreman's hook, in that the shaft is typically longer. It is used in hay bucking on farms to secure and move bales of hay, which are otherwise awkward to pick up manually.
Gardening
In gardening and agriculture, a variant with a long shaft is used to move large plants. A hook is placed in either side of the baled roots, allowing workers to carry or place the heavy load.
Forestry
Called a "Packhaken", "Hebehaken", or "Forsthaken" in German, this type is used in forestry mainly to lift or move firewood. In Sweden, this tool, though slightly different, is called a "timmerkrok", which translates as "timberhook". It is used mainly by two people to move logs by hooking them in each end.
| Technology | Hand tools | null |
1133088 | https://en.wikipedia.org/wiki/Elastic%20scattering | Elastic scattering | Elastic scattering is a form of particle scattering in scattering theory, nuclear physics and particle physics. In this process, the internal states of the particles involved stay the same. In the non-relativistic case, where the relative velocities of the particles are much less than the speed of light, elastic scattering simply means that the total kinetic energy of the system is conserved. At relativistic velocities, elastic scattering also requires the final state to have the same number of particles as the initial state and for them to be of the same kind.
Rutherford scattering
When the incident particle, such as an alpha particle or electron, is diffracted in the Coulomb potential of atoms and molecules, the elastic scattering process is called Rutherford scattering. In many electron diffraction techniques like reflection high energy electron diffraction (RHEED), transmission electron diffraction (TED), and gas electron diffraction (GED), where the incident electrons have sufficiently high energy (>10 keV), the elastic electron scattering becomes the main component of the scattering process and the scattering intensity is expressed as a function of the momentum transfer defined as the difference between the momentum vector of the incident electron and that of the scattered electron.
Optical elastic scattering
In Thomson scattering light interacts with electrons (this is the low-energy limit of Compton scattering).
In Rayleigh scattering a medium composed of particles whose sizes are much smaller than the wavelength scatters light sideways. In this scattering process, the energy (and therefore the wavelength) of the incident light is conserved and only its direction is changed. In this case, the scattering intensity is inversely proportional to the fourth power of the reciprocal wavelength of the light.
Nuclear particle physics
For particles with the mass of a proton or greater, elastic scattering is one of the main methods by which the particles interact with matter. At relativistic energies, protons, neutrons, helium ions, and HZE ions will undergo numerous elastic collisions before they are dissipated. This is a major concern with many types of ionizing radiation, including galactic cosmic rays, solar proton events, free neutrons in nuclear weapon design and nuclear reactor design, spaceship design, and the study of the Earth's magnetic field. In designing an effective biological shield, proper attention must be made to the linear energy transfer of the particles as they propagate through the shield. In nuclear reactors, the neutron's mean free path is critical as it undergoes elastic scattering on its way to becoming a slow-moving thermal neutron.
Besides elastic scattering, charged particles also undergo effects from their elementary charge, which repels them away from nuclei and causes their path to be curved inside an electric field. Particles can also undergo inelastic scattering and capture due to nuclear reactions. Protons and neutrons do this more often than heavier particles. Neutrons are also capable of causing fission in an incident nucleus. Light nuclei like deuterium and lithium can combine in nuclear fusion.
| Physical sciences | Particle physics: General | Physics |
1135529 | https://en.wikipedia.org/wiki/Hardiness%20%28plants%29 | Hardiness (plants) | Hardiness of plants describes their ability to survive adverse growing conditions. It is usually limited to discussions of climatic adversity. Thus a plant's ability to tolerate cold, heat, drought, flooding, or wind are typically considered measurements of hardiness. Hardiness of plants is defined by their native extent's geographic location: longitude, latitude and elevation. These attributes are often simplified to a hardiness zone. In temperate latitudes, the term most often describes resistance to cold, or "cold-hardiness", and is generally measured by the lowest temperature a plant can withstand.
Hardiness of a plant may be divided into two categories: tender, and hardy. Tender plants are those killed by freezing temperatures, while hardy plants survive freezing—at least down to certain temperatures, depending on the plant. "Half-hardy" is a term used sometimes in horticulture to describe bedding plants which are sown in heat in winter or early spring, and planted outside after all danger of frost has passed. "Fully hardy" usually refers to plants being classified under the Royal Horticultural Society classifications, and can often cause confusion to those not using this method. When this distinction is made a fully tropical plant that requires hot temperatures to grow and display is termed a "tender" plant.
Plants vary greatly in their tolerance of growing conditions, and are capable of adaptation to changes in climate on their own to some extent. The selective breeding of varieties capable of withstanding particular climates forms an important part of agriculture and horticulture. Part of the work of nursery growers of plants consists of cold hardening, or hardening off their plants, to prepare them for likely conditions in later life.
Winter hardiness
Winter-hardy plants grow during the winter, or at least remain healthy and dormant. Apart from hardy evergreens, these include many cultivated plants, including varieties of cabbage and broccoli, and all kinds of carrot. Some bulbs – such as tulips – need cold winters to bloom, while others – such as freesia – can survive a freezing winter. Many domestic plants are assigned a hardiness zone that specifies the climates in which they can survive. Winter gardens are dependent upon the cultivation of winter-hardy plants.
Mechanisms
Woody plants survive freezing temperatures by suppressing the formation of ice in living cells or by allowing water to freeze in plant parts that are not affected by ice formation. The common mechanism for woody plants to survive down to –40 °C (–40 °F) is supercooling. Woody plants that survive lower temperatures are dehydrating their cells, allowing water to freeze between cell walls and the cells to survive. Plants normally considered hardy may not survive freezing if they are not readily acclimated, which renders them unable to use these mechanisms. Some plants are able to withstand freezing temperatures year-round, like those living in Antarctica such as Colobanthus quitensis.
Hardiness ratings
Various hardiness ratings are published. In the United States (US), the most widely used is the U.S. Department of Agriculture (USDA) system of hardiness zones based on average minimum yearly temperatures. This system was developed specifically for the extremely diverse range of conditions in the US, from baking desert to frozen tundra. Another commonly used system is the Sunset Climate Zone system. This system is much more specific to climates (i.e. precipitation, temperature, and humidity based) and less dependent on the yearly minimum.
In contrast the United Kingdom (UK) and Western Europe have an oceanic climate, and experience a narrower range of temperatures that is tempered by the presence of the Gulf Stream. This results in areas like western Scotland experiencing relatively mild winter conditions and cool summers that enable the growing of some subtropical plants, despite being well to the north of subtropical climate areas. The Royal Horticultural Society has published a set of hardiness ratings applicable to the UK. The ratings range from H1a to H7. H1a, higher than , applies to tropical plants permanently under glass in heat; while H7, below , applies to very cold-tolerant plants such as heathers. Most outdoor plants in the UK fall within the range H4, (hardy in the average winter) to H5, (hardy in a cold winter). Also, the average minimum temperature in the UK is much warmer than the average minimums in much of the northern US, while the average maximums in the northern US are often much warmer in summer than the UK.
In addition to cold tolerance, plant hardiness has been observed to be linked to how much stress specific plants are undergoing into the winter, or even how fast the onset of cold weather is in a specific year. This means that often stressed plants will exhibit less cold tolerance than plants that have been well maintained. Plants may also die if the winter changes from balmy to exceptionally cold in a short period of time.
Hardy plants
Hardy plants are generally known as plants that can tolerate a variety of extreme (both cold and hot) temperatures , harsh environments, and poor or eroded soils, such as Milkweed and Cranesbill; and economically advantageous because of their little or negligible maintenance .
Other hardiness
Tolerance to other climatic factors are often omitted from the "hardiness" discussion. This section provides a very brief overview of the topics.
Heat tolerance may be displayed alongside cold-hardiness, such as in the American Horticultural Society's "Heat Zones". See and Breeding for heat stress tolerance.
For drought tolerance, see and Breeding for drought stress tolerance.
For flood tolerance, see cited review.
The mechanisms of wind tolerance are still poorly understood. Wind not only causes mechanical stress, but also makes the plant colder and dryer. Plants may evolutionally adapt to wind by becoming smaller (to avoid wind) or bigger (to resist the wind). Physiologically, the mechanical force also causes a cascade of signals in the living plant.
| Technology | Horticulture | null |
1136447 | https://en.wikipedia.org/wiki/Microcanonical%20ensemble | Microcanonical ensemble | In statistical mechanics, the microcanonical ensemble is a statistical ensemble that represents the possible states of a mechanical system whose total energy is exactly specified. The system is assumed to be isolated in the sense that it cannot exchange energy or particles with its environment, so that (by conservation of energy) the energy of the system does not change with time.
The primary macroscopic variables of the microcanonical ensemble are the total number of particles in the system (symbol: ), the system's volume (symbol: ), as well as the total energy in the system (symbol: ). Each of these is assumed to be constant in the ensemble. For this reason, the microcanonical ensemble is sometimes called the ensemble.
In simple terms, the microcanonical ensemble is defined by assigning an equal probability to every microstate whose energy falls within a range centered at . All other microstates are given a probability of zero. Since the probabilities must add up to 1, the probability is the inverse of the number of microstates within the range of energy,
The range of energy is then reduced in width until it is infinitesimally narrow, still centered at . In the limit of this process, the microcanonical ensemble is obtained.
Applicability
Because of its connection with the elementary assumptions of equilibrium statistical mechanics (particularly the postulate of a priori equal probabilities), the microcanonical ensemble is an important conceptual building block in the theory. It is sometimes considered to be the fundamental distribution of equilibrium statistical mechanics. It is also useful in some numerical applications, such as molecular dynamics. On the other hand, most nontrivial systems are mathematically cumbersome to describe in the microcanonical ensemble, and there are also ambiguities regarding the definitions of entropy and temperature. For these reasons, other ensembles are often preferred for theoretical calculations.
The applicability of the microcanonical ensemble to real-world systems depends on the importance of energy fluctuations, which may result from interactions between the system and its environment as well as uncontrolled factors in preparing the system. Generally, fluctuations are negligible if a system is macroscopically large, or if it is manufactured with precisely known energy and thereafter maintained in near isolation from its environment. In such cases the microcanonical ensemble is applicable. Otherwise, different ensembles are more appropriate – such as the canonical ensemble (fluctuating energy) or the grand canonical ensemble (fluctuating energy and particle number).
Properties
Thermodynamic quantities
The fundamental thermodynamic potential of the microcanonical ensemble is entropy. There are at least three possible definitions, each given in terms of the phase volume function . In classical mechanics this is the volume of the region of phase space where the energy is less than . In quantum mechanics is roughly the number of energy eigenstates with energy less than ; however this must be smoothed so that we can take its derivative (see the Precise expressions section for details on how this is done). The definitions of microcanonical entropy are:
In the microcanonical ensemble, the temperature is a derived quantity rather than an external control parameter. It is defined as the derivative of the chosen entropy with respect to energy. For example, one can define the "temperatures" and as follows:
Like entropy, there are multiple ways to understand temperature in the microcanonical ensemble. More generally, the correspondence between these ensemble-based definitions and their thermodynamic counterparts is not perfect, particularly for finite systems.
The microcanonical pressure and chemical potential are given by:
Phase transitions
Under their strict definition, phase transitions correspond to nonanalytic behavior in the thermodynamic potential or its derivatives. Using this definition, phase transitions in the microcanonical ensemble can occur in systems of any size. This contrasts with the canonical and grand canonical ensembles, for which phase transitions can occur only in the thermodynamic limit – i.e., in systems with infinitely many degrees of freedom. Roughly speaking, the reservoirs defining the canonical or grand canonical ensembles introduce fluctuations that "smooth out" any nonanalytic behavior in the free energy of finite systems. This smoothing effect is usually negligible in macroscopic systems, which are sufficiently large that the free energy can approximate nonanalytic behavior exceedingly well. However, the technical difference in ensembles may be important in the theoretical analysis of small systems.
Information entropy
For a given mechanical system (fixed , ) and a given range of energy, the uniform distribution of probability over microstates (as in the microcanonical ensemble) maximizes the ensemble average .
Thermodynamic analogies
Early work in statistical mechanics by Ludwig Boltzmann led to his eponymous entropy equation for a system of a given total energy, , where is the number of distinct states accessible by the system at that energy. Boltzmann did not elaborate too deeply on what exactly constitutes the set of distinct states of a system, besides the special case of an ideal gas. This topic was investigated to completion by Josiah Willard Gibbs who developed the generalized statistical mechanics for arbitrary mechanical systems, and defined the microcanonical ensemble described in this article. Gibbs investigated carefully the analogies between the microcanonical ensemble and thermodynamics, especially how they break down in the case of systems of few degrees of freedom. He introduced two further definitions of microcanonical entropy that do not depend on – the volume and surface entropy described above. (Note that the surface entropy differs from the Boltzmann entropy only by an -dependent offset.)
The volume entropy and associated temperature are closely analogous to thermodynamic entropy and temperature. It is possible to show exactly that
( is the ensemble average pressure) as expected for the first law of thermodynamics. A similar equation can be found for the surface entropy (or Boltzmann entropy ) and its associated temperature , however the "pressure" in this equation is a complicated quantity unrelated to the average pressure.
The microcanonical temperatures and are not entirely satisfactory in their analogy to temperature as defined using a canonical ensemble. Outside of the thermodynamic limit, a number of artefacts occur.
Nontrivial result of combining two systems: Two systems, each described by an independent microcanonical ensemble, can be brought into thermal contact and be allowed to equilibriate into a combined system also described by a microcanonical ensemble. Unfortunately, the energy flow between the two systems cannot be predicted based on the initial s. Even when the initial s are equal, there may be energy transferred. Moreover, the of the combination is different from the initial values. This contradicts the intuition that temperature should be an intensive quantity, and that two equal-temperature systems should be unaffected by being brought into thermal contact.
Strange behavior for few-particle systems: Many results such as the microcanonical Equipartition theorem acquire a one- or two-degree of freedom offset when written in terms of . For a small systems this offset is significant, and so if we make the analogue of entropy, several exceptions need to be made for systems with only one or two degrees of freedom.
Spurious negative temperatures: A negative occurs whenever the density of states is decreasing with energy. In some systems the density of states is not monotonic in energy, and so can change sign multiple times as the energy is increased.
The preferred solution to these problems is avoid use of the microcanonical ensemble. In many realistic cases a system is thermostatted to a heat bath so that the energy is not precisely known. Then, a more accurate description is the canonical ensemble or grand canonical ensemble, both of which have complete correspondence to thermodynamics.
Precise expressions for the ensemble
The precise mathematical expression for a statistical ensemble depends on the kind of mechanics under consideration – quantum or classical – since the notion of a "microstate" is considerably different in these two cases. In quantum mechanics, diagonalization provides a discrete set of microstates with specific energies. The classical mechanical case involves instead an integral over canonical phase space, and the size of microstates in phase space can be chosen somewhat arbitrarily.
To construct the microcanonical ensemble, it is necessary in both types of mechanics to first specify a range of energy. In the expressions below the function (a function of , peaking at with width ) will be used to represent the range of energy in which to include states. An example of this function would be
or, more smoothly,
Quantum mechanical
A statistical ensemble in quantum mechanics is represented by a density matrix, denoted by . The microcanonical ensemble can be written using bra–ket notation, in terms of the system's energy eigenstates and energy eigenvalues. Given a complete basis of energy eigenstates , indexed by , the microcanonical ensemble is
where the are the energy eigenvalues determined by (here is the system's total energy operator, i. e., Hamiltonian operator). The value of is determined by demanding that is a normalized density matrix, and so
The state volume function (used to calculate entropy) is given by
The microcanonical ensemble is defined by taking the limit of the density matrix as the energy width goes to zero, however a problematic situation occurs once the energy width becomes smaller than the spacing between energy levels. For very small energy width, the ensemble does not exist at all for most values of , since no states fall within the range. When the ensemble does exist, it typically only contains one (or two) states, since in a complex system the energy levels are only ever equal by accident (see random matrix theory for more discussion on this point). Moreover, the state-volume function also increases only in discrete increments, and so its derivative is only ever infinite or zero, making it difficult to define the density of states. This problem can be solved by not taking the energy range completely to zero and smoothing the state-volume function, however this makes the definition of the ensemble more complicated, since it becomes then necessary to specify the energy range in addition to other variables (together, an ensemble).
Classical mechanical
In classical mechanics, an ensemble is represented by a joint probability density function defined over the system's phase space. The phase space has generalized coordinates called , and associated canonical momenta called .
The probability density function for the microcanonical ensemble is:
where
is the total energy (Hamiltonian) of the system, a function of the phase ,
is an arbitrary but predetermined constant with the units of , setting the extent of one microstate and providing correct dimensions to .
is an overcounting correction factor, often used for particle systems where identical particles are able to change place with each other.
Again, the value of is determined by demanding that is a normalized probability density function:
This integral is taken over the entire phase space. The state volume function (used to calculate entropy) is defined by
As the energy width is taken to zero, the value of decreases in proportion to as .
Based on the above definition, the microcanonical ensemble can be visualized as an infinitesimally thin shell in phase space, centered on a constant-energy surface. Although the microcanonical ensemble is confined to this surface, it is not necessarily uniformly distributed over that surface: if the gradient of energy in phase space varies, then the microcanonical ensemble is "thicker" (more concentrated) in some parts of the surface than others. This feature is an unavoidable consequence of requiring that the microcanonical ensemble is a steady-state ensemble.
Examples
Ideal gas
The fundamental quantity in the microcanonical ensemble is , which is equal to the phase space volume compatible with given . From , all thermodynamic quantities can be calculated. For an ideal gas, the energy is independent of the particle positions, which therefore contribute a factor of to . The momenta, by contrast, are constrained to a -dimensional (hyper-)spherical shell of radius ; their contribution is equal to the surface volume of this shell. The resulting expression for is:
where is the gamma function, and the factor has been included to account for the indistinguishability of particles (see Gibbs paradox). In the large limit, the Boltzmann entropy is
This is also known as the Sackur–Tetrode equation.
The temperature is given by
which agrees with analogous result from the kinetic theory of gases. Calculating the pressure gives the ideal gas law:
Finally, the chemical potential is
Ideal gas in a uniform gravitational field
The microcanonical phase volume can also be calculated explicitly for an ideal gas in a uniform gravitational field.
The results are stated below for a 3-dimensional ideal gas of particles, each with mass , confined in a thermally isolated container that is infinitely long in the z-direction and has constant cross-sectional area . The gravitational field is assumed to act in the minus z direction with strength . The phase volume is
where is the total energy, kinetic plus gravitational.
The gas density as a function of height can be obtained by integrating over the phase volume coordinates. The result is:
Similarly, the distribution of the velocity magnitude (averaged over all heights) is
The analogues of these equations in the canonical ensemble are the barometric formula and the Maxwell–Boltzmann distribution, respectively. In the limit , the microcanonical and canonical expressions coincide; however, they differ for finite . In particular, in the microcanonical ensemble, the positions and velocities are not statistically independent. As a result, the kinetic temperature, defined as the average kinetic energy in a given volume , is nonuniform throughout the container:
By contrast, the temperature is uniform in the canonical ensemble, for any .
| Physical sciences | Statistical mechanics | Physics |
1137144 | https://en.wikipedia.org/wiki/Antisense%20therapy | Antisense therapy | Antisense therapy is a form of treatment that uses antisense oligonucleotides (ASOs) to target messenger RNA (mRNA). ASOs are capable of altering mRNA expression through a variety of mechanisms, including ribonuclease H mediated decay of the pre-mRNA, direct steric blockage, and exon content modulation through splicing site binding on pre-mRNA. Several ASOs have been approved in the United States, the European Union, and elsewhere.
Nomenclature
The common stem for antisense oligonucleotides drugs is -rsen. The substem -virsen designates antiviral antisense oligonucleotides.
Antisense Oligonucleotide Development
Developments in ASO modification are separated into three generations. Generation one is called backbone-modified and focuses on the phosphodiester group of the nucleotide. This impacts inter-nucleotide binding. These modifications led to better distribution, reduced urinary excretion, and prolonged residence time of the ASOs in the cell. Some examples of first generation modifications include the addition of a phosphorothioate group (PS), methyl group, or nitrogen. The most common is the phosphorothioate group (PS) in which the oxygen atoms of a phosphodiester group are replaced with sulfur atoms, greatly improving efficacy and reducing degradation. Generation two is sugar-modified, focused on the ribose sugar of the nucleotide. This generation saw improved binding affinity while reducing degradation. Some examples of generation two modifications are the substitution of R group with morpholine group (MO) and the usage of phosphorodiamidate morpholino oligomer (PMO) and thiomorpholine oligomer (TMO) as linkages between the ribose sugar and phosphodiester group in the backbone. Generation three is nucleobase-modified, the least common type of modification. These modifications enhanced binding affinity and cell penetration while reducing degradation and off-target effects. Examples include the introduction of G-clamps, pseudoisocytosine, and the substitution of bases with amine, thione, halogen, alkyl, alkenyl, or alkynyl groups.
Pharmacokinetics and pharmacodynamics
Half-life and stability
ASO-based drugs employ highly modified, single-stranded chains of synthetic nucleic acids that achieve wide tissue distribution with very long half-lives. For instance, many ASO-based drugs contain phosphorothioate substitutions and 2' sugar modifications to inhibit nuclease degradation enabling vehicle-free delivery to cells.
In vivo delivery
Phosphorothioate ASOs can be delivered to cells without the need of a delivery vehicle. ASOs do not penetrate the blood brain barrier when delivered systemically but they can distribute across the neuraxis if injected in the cerebrospinal fluid typically by intrathecal administration. Newer formulations using conjugated ligands greatly enhances delivery efficiency and cell-type specific targeting.
Approved therapies
Amyotrophic lateral sclerosis
Tofersen (marketed as Qalsody) was approved by the FDA for the treatment of SOD1- associated amyotrophic lateral sclerosis (ALS) in 2023. It was developed by Biogen under a licensing agreement with Ionis Pharmaceuticals. In trials the drug was found to lower levels of an ALS biomarker, neurofilament light change, and in long-term trial extensions to slow disease. Under the terms of the FDA's accelerated approval program, a confirmatory study will be conducted in presymptomatic gene carriers to provide additional evidence.
Batten disease
Milasen is a novel individualized therapeutic agent that was designed and approved by the FDA for the treatment of Batten disease. This therapy serves as an example of personalized medicine.
In 2019, a report was published detailing the development of milasen, an antisense oligonucleotide drug for Batten disease, under an expanded-access investigational clinical protocol authorized by the Food and Drug Administration (FDA). Milasen "itself remains an investigational drug, and it is not suited for the treatment of other patients with Batten's disease" because it was customized for a single patient's specific mutation. However it is an example of individualized genomic medicine therapeutical intervention.
Cytomegalovirus retinitis
Fomivirsen (marketed as Vitravene), was approved by the U.S. FDA in August 1998, as a treatment for cytomegalovirus retinitis.
Duchenne muscular dystrophy
Several morpholino oligos have been approved to treat specific groups of mutations causing Duchenne muscular dystrophy. In September 2016, eteplirsen (ExonDys51) received FDA approval for the treatment of cases that can benefit from skipping exon 51 of the dystrophin transcript. In December 2019, golodirsen (Vyondys 53) received FDA approval for the treatment of cases that can benefit from skipping exon 53 of the dystrophin transcript. In August 2020, viltolarsen (Viltepso) received FDA approval for the treatment of cases that can benefit from skipping exon 53 of the dystrophin transcript.
Familial chylomicronaemia syndrome
Volanesorsen was approved by the European Medicines Agency (EMA) for the treatment of familial chylomicronaemia syndrome in May 2019.
Familial hypercholesterolemia
In January 2013 mipomersen (marketed as Kynamro) was approved by the FDA for the treatment of homozygous familial hypercholesterolemia.
Hereditary transthyretin-mediated amyloidosis
Inotersen received FDA approval for the treatment of hereditary transthyretin-mediated amyloidosis in October 2018. The application for inotersen was granted orphan drug designation. It was developed by Ionis Pharmaceuticals and licensed to Akcea Therapeutics. Patisiran (sold under Onpattro) was developed by Alnylam Pharmaceuticals, and also approved for use in the US and EU in 2018 with orphan drug designation. Its mechanism-of-action is the active substance of small interfering RNA (siRNA), which allows it to interfere with and block the production of trasnthyretin. As such, it was the first FDA-approved siRNA therapeutic.
Spinal muscular atrophy
In 2004, development of an antisense therapy for spinal muscular atrophy began. Over the following years, an antisense oligonucleotide later named nusinersen was developed by Ionis Pharmaceuticals under a licensing agreement with Biogen. In December 2016, nusinersen received regulatory approval from FDA and soon after, from other regulatory agencies worldwide.
Investigational therapies
Current clinical trials
As of 2020 more than 50 antisense oligonucleotides were in clinical trials, including over 25 in advanced clinical trials (phase II or III).
Phase III trials
Hereditary transthyretin-mediated amyloidosis
A follow-on drug to Inotersen is being developed by Ionis Pharmaceuticals and under license to Akcea Therapeutics for hereditary transthyretin-mediated amyloidosis. In this formulation the ASO is conjugated to N-Acetylgalactosamine enabling hepatocyte-specific delivery, greatly reducing dose requirements and side effect profile while increasing the level of transthyretin reduction in patients.
Huntington's disease
Tominersen (also known as IONIS-HTTRx and RG6042) was tested in a phase 3 trial for Huntington's disease although this trial was discontinued on March 21, 2021, due to lack of efficacy. It is currently licensed to Roche by Ionis Pharmaceuticals.
Phase I and II trials
Clinical trials are ongoing for several diseases and conditions including:
Acromegaly, age related macular degeneration, Alzheimer's disease, amyotrophic lateral sclerosis, autosomal dominant retinitis pigmentosa, beta thalassemia, cardiovascular disease, elevated level of lipoprotein(a), centronuclear myopathy, coagulopathies, cystic fibrosis, Duchenne muscular dystrophy, diabetes, epidermolysis bullosa dystrophica, familial chylomicronemia syndrome, frontotemporal dementia, Fuchs' dystrophy, hepatitis B, hereditary angioedema, hypertension, IgA nephropathy, Leber's hereditary optic neuropathy, multiple system atrophy, non-alcoholic fatty liver disease, Parkinson's disease, prostate cancer, Stargardt disease, STAT3-expressing cancers, Usher syndrome.
Preclinical development
Several ASOs are currently being investigated in disease models for Alexander disease, ATXN2 (gene) and FUS (gene) amyotrophic lateral sclerosis, Angelman syndrome, Lafora disease, lymphoma, multiple myeloma, myotonic dystrophy, Parkinson's disease, Pelizaeus–Merzbacher disease, and prion disease, Rett syndrome, spinocerebellar Ataxia Type 3.
| Technology | Biotechnology | null |
1137210 | https://en.wikipedia.org/wiki/Pocketknife | Pocketknife | A pocketknife is a knife with one or more blades that fold into the handle. They are also known as jackknives, folding knives, EDC knife, or may be referred to as a penknife, though a penknife may also be a specific kind of pocketknife. A typical blade length is .
Pocketknives are versatile tools, and may be used for anything from whittling and woodcarving, to butchering small game, gutting and filleting small fish, aiding in the preparation of tinder and kindling for fires, boring holes in soft material, to opening an envelope, cutting twine, slicing fruits and vegetables or as a means of self-defense.
Pocketknives may also be used in conjunction with other tools and equipment for woodcraft and bushcraft.
Specialised designs are also used for mushroom hunting and gardening. Pocketknives designed for gardening include pruning knives, which are folding knives with long curved blades used for pruning, trimming cuttings, taking buds and preparing material for grafting.
History
The earliest known pocketknives date to at least the early Iron Age. A pocketknife with a bone handle was found at the Hallstatt Culture type site in Austria, dating to around 600–500 BCE. Iberian folding-blade knives made by indigenous artisans and craftsmen and dating to the pre-Roman era have been found in Spain. Many folding knives from the Viking Age have been found. They carried some friction binders, but more often they seem to have used folding knives that used a closure to keep the blade open.
Peasant knife
The peasant knife, farmer knife, or penny knife is the original and most basic design of a folding pocketknife, using a simple pivoted blade that folds in and out of the handle freely, without a backspring, slipjoint, or blade locking mechanism. The first peasant knives date to the pre-Roman era, but were not widely distributed nor affordable by most people until the advent of limited production of such knives in cutlery centers such as Sheffield, England commencing around 1650, with large-scale production starting around the year 1700 with models such as Fuller's Penny Knife and the Wharncliffe Knife. Some peasant knives used a bolster or tensioning screw at the blade to apply friction to the blade tang in order to keep the blade in the open position. The smallest (Nos. 2–5) Opinel knives are an example of the peasant knife. The knife's low cost made it a favorite of small farmers, herdsmen, and gardeners in Europe and the Americas during the late 19th and early 20th century.
Slip joint knife
Most pocketknives for light duty are slipjoints. This means that the blade does not lock but, once opened, is held in place by tension from a flat bar or leaf-type backspring that allows the blade to fold if a certain amount of pressure is applied. The first spring-back knives were developed around 1660 in England, but were not widely available until the Industrial Revolution and development of machinery capable of mass production. Many locking knives have only one blade that is as large as can be fitted into the handle, because the locking mechanism relies on a spring-loaded latch built into the spine or frame of the handle to lock it and it is difficult to build in multiple levers, one for each blade. Slipjoints tend to be smaller than other typical pocketknives.
Some popular patterns of slipjoint knives include:
Multi-tool knives
Multi-tool knives formerly consisted of variations on the American camper style or the Swiss Army knives manufactured by Victorinox and Wenger. However, the concept of a multitool knife has undergone a revolution thanks in part to an avalanche of new styles, sizes, and tool presentation concepts. These new varieties often incorporate a pair of pliers and other tools in conjunction with one or more knife blade styles, either locking or nonlocking.
Multitool knives often have more than one blade, including an assortment of knife blade edges (serrated, plain, saws) as well as a selection of other tools such as bottle openers, corkscrews, and scissors. A large tool selection is the signature of the Swiss Army Knife. These knives are produced by Victorinox and Wenger and issued to military services and sold to the public. Similar to the Swiss Army knife is the German Army knife, with two blades opening from each side and featuring hard plastic grips and aluminum liners. The U.S. Military utility knife (MIL-K-818), issued by the United States Army, Navy, and Marine Corps, was made for many years by the Camillus Cutlery Company and Imperial Schrade as well as many other companies. It was originally produced with carbon steel blades and brass liners (both vulnerable to corrosion), but with the onset of the Vietnam War was modified to incorporate all-stainless steel construction. The current-issue U.S. military utility knife has textured stainless grips and four stainless blades/tools opening on both sides in the camper or scout pattern and has an extremely large clevis or bail. The Victorinox Swiss Army Soldier Knife has been issued a National Stock Number (NSN) to be authorized for issue to US service members. This is the first time a US military issued utility/pocket knife that was not manufactured in the United States has been adopted.
Miscellaneous designs
Another style of folding, non-locking knife is the friction-folder. These use simple friction between the blade and scales to hold the blade in place once opened; an example is the Japanese higonokami.
An electrician's knife typically has a locking screwdriver blade but a non-locking knife blade. The two-blade Camillus Electrician's knife (the US military version is known as a TL-29) was the inspiration for the development of the linerlock.
A credit card knife is usually a very thin knife that is the shape and size of a credit card, either when folded into a knife shape or unfolded for storage. It is designed to be carried in a wallet along with regular credit cards. Some credit card knives can contain other small tools, such as tweezers, or toothpicks.
A ballpoint pen knife is generally a pen with a concealed knife inside, which can be used as a letter opener or as a self-defense weapon.
Gravity lock ball bearing pocketknife, locks with sphere dropping in and out into grove. Has to be oriented tip up and pressed to release, tip down and pressed to lock.
Lock-blade knives
Knives with locking blades, often referred to as lock-blade knives or clasp knives, have a locking mechanism that locks the blade into its fully opened position. This lock must be released in a distinct action before the knife can be folded. The lock-blade knife improves safety by preventing accidental blade closure while cutting. It is this locking blade feature that differentiates the lock-blade knife from either the peasant knife or the slipjoint spring-back knife. Locking knives also tend to be larger: it is easier to fit a lock into a larger frame, and larger knives are more likely to be used for more forceful kinds of work. The cost of a locking mechanism is also proportionally less than it would be on a smaller, and generally cheaper, knife.
Lock-blade knives have been dated to the 15th century. In Spain, one early lock-blade design was the Andalusian clasp knife popularly referred to as the navaja. Opinel knives use a twist lock, consisting of a metal ferrule or barrel ring that is rotated to lock the blade either open or closed. In the late 20th century lock-blade pocketknives were popularized and marketed on a wider scale. Companies such as Buck Knives, Camillus, Case, and Gerber, created a wide range of products with locks of various types. The most popular form, the lockback knife, was popularized by Buck Knives in the 1960s, so much that the eponymous term "buck knife" was used to refer to lockback knives that were not manufactured by Buck.
The lockback's blade locking mechanism is a refinement of the slipjoint design; both use a strong backspring located along the back of the knife handle. However, the lockback design incorporates a hook or lug on the backspring, which snaps into a corresponding notch on the blade's heel when the blade is fully opened, locking the blade into position. Closing the blade requires the user to apply pressure to the spring-loaded bar located towards the rear of the knife handle to disengage the hook from the notch and thus release the blade.
The Walker Linerlock, invented by knifemaker Michael Walker, and the framelock came to prominence in the 1980s. In both designs the liner inside the knife is spring-loaded to engage the rear of the blade when open and thus hold it in place. In the case of the framelock, the liner is the handle, itself. The Swiss Army knife product range has adopted dual linerlocks on their 111 mm models. Some models feature additional "positive" locks, which essentially ensure that the blade cannot close accidentally. CRKT has patented an "Auto-LAWKS" device, which features a second sliding switch on the hilt. It can operate as any linerlock knife if so desired, but if the user slides the second control up after opening, it places a wedge between the linerlock and the frame, preventing the lock from disengaging until the second device is disabled.
Tactical folding knife
Buck's original lockback knife was originally marketed as a "folding hunting knife" and while it became popular with sportsmen, it saw use with military personnel as it could perform a variety of tasks. Custom knife makers began making similar knives, in particular was knifemaker Bob Terzuola. Terzuola is credited with coining the phrase "Tactical Folder".
In the early 1990s, tactical folding knives became popular in the U.S.A. The trend began with custom knifemakers such as Bob Terzuola, Michael Walker, Allen Elishewitz, Mel Pardue, Ernest Emerson, Ken Onion, Chris Reeve, Rick Hinderer, Warren Thomas, and Warren Osbourne. These knives were most commonly built as linerlocks. Blade lengths varied from , but the most typical models never exceeded in blade length for legal reasons in most US jurisdictions.
In response to the demand for these knives, production companies offered mass-produced tactical folding knives. Companies such as Benchmade, Kershaw Knives, Buck Knives, Gerber, CRKT, Spyderco and Cold Steel collaborated with tactical knifemakers; in some cases retaining them as full-time designers. Tactical knifemakers such as Ernest Emerson and Chris Reeve went so far as to open their own mass-production factories.
Presenting any folding knife as a weapon, rather than a utility tool, has met criticism. Many who've studied knife fighting point out that even the strongest locking mechanisms have some risk of failure, so a folding knife is never going to be as reliable as a fixed-blade combat knife. Lynn Thompson, martial artist and CEO of Cold Steel, noted in an article in Black Belt magazine that most tactical folding knives are too short to be of much use in a knife fight. And while his company does make and sell a tactical folder, it is not the best option during an actual fight.
The case for tactical folding knives, then, is that a less-than-ideal knife in your pocket is always more useful than an ideal knife left at home. A 10-inch fixed-blade Bowie knife, for example, may be far better for combat, but it is far less practical — and often illegal — to carry around in day-to-day life. And should there be a need to take defensive action, there's rarely enough forewarning to plan ahead which knife to carry that day.
Other features
Traditional folding knives are opened using nail-nicks, or slots where the user's fingernail would enter to pull the blade out of the handle. This became somewhat cumbersome and required use of two hands, so there were innovations to remedy that. The thumb-stud, a small stud on the blade that allows for one-handed opening, led the way for more innovations. One of these is the thumb hole: a Spyderco patent where the user presses the pad of the thumb against a hole and opens the blade by rotating the thumb similarly to using the thumb-stud.
Another innovation of Sal Glesser, Spyderco founder, was the clip system, which he named a "Clip-it". Clips are usually metal or plastic and similar to the clips found on pens except thicker. Clips allow the knife to be easily accessible, while keeping it lint-free and unscathed by pocket items such as coins. Assisted opening systems have been pioneered by makers like Ken Onion with his "Speed-Safe" mechanism and Ernest Emerson's Wave system, where a hook catches the user's pocket upon removal and the blade is opened during a draw.
One of the first one handed devices was the automatic spring release, also known as a switchblade. An innovation to pocketknives made possible by the thumb-stud is the replaceable blade insert developed in 1999 by Steven Overholt (U.S. Patent no. 6,574,868), originally marketed by TigerSharp Technologies and as of 2007 by Clauss. Some systems are somewhat between assisted opening and the normal thumb stud. CRKT knives designed by Harold "Kit" Carson often incorporate a "Carson Flipper", which is a small protrusion on the rear of the base of the blade such that it protrudes out the obverse side of the handle (when closed). By using an index finger and a very slight snapping of the wrist, the knife opens very quickly, appearing to operate like a spring assisted knife. When opened, the protrusion is between the base of the sharp blade and the user's index finger, preventing any accidental slipping of the hand onto the blade. Some designs feature a second "Flipper" on the opposite side of the blade, forming a small "hilt guard" such as a fixed blade knife has, which can prevent another blade from sliding up into the hilt in combat. These "flippers" are now being found on other brands of knives as well, such as Kershaw, even cheaper knives, including certain versions of Schrade's Snowblind tactical folders, and numerous others.
Legal issues
Pocketknives are legal to own in most countries, but may face legal restrictions on their use. While pocketknives are almost always designed as tools, they do have the potential to be considered by legal authorities as weapons.
In the United States, knives are regulated by federal, state, and municipal laws. Some jurisdictions prohibit the possession or use of pocketknives that feature locking blades. Others prohibit certain blade styles perceived by law enforcement and legal authorities as optimal for offensive fighting, transforming the pocketknife from a utility tool into a deadly weapon. These might include knives with dirk, dagger (double-edge), bowie, or stiletto blades. In some jurisdictions it is illegal to conceal knives larger than a certain size or with blades over a certain length, particularly when combined with locking blade mechanisms. The possession or carrying of a folding pocketknife with a quick-opening mechanism such as a gravity knife, butterfly knife (balisong), or switchblade may be prohibited. Under U.S. federal law, switchblades and ballistic knives are banned from interstate shipment, sale, or import, or possession on federal or Indian lands or U.S. possessions and may be prohibited entirely in some states. Knives of any size or configuration may be prohibited by federal or state laws in certain designated areas or places, such as schools, courthouses, jails, power plants, or airports.
In the United Kingdom it is illegal to carry a folding knife having a blade with a cutting edge of more than 3 inches (just over 7.6 cm) in length in public without "good reason". The terms "in public" and "good reason" are not defined, but examples of "religious duty", "national dress" and "requirement of employment or hobby" are given. It is left up to a police officer's individual subjective discretion, and ultimately a magistrate to decide if a knife is being carried "in public", and for a "good reason". Folding knives with blades of or less may be carried without needing to provide "good reason" so long as the blade is not capable of being locked in the open position. However, it is illegal to have the intention of using any object in public as a weapon, meaning that even a knife that is legal to carry without needing "good reason" may still be found to be illegal if the police officer has grounds to suspect it will be used as a weapon. The onus lies on the officer to prove that intent. Recent court decisions in the U.K. have made it easier for public prosecutors to obtain knife possession convictions by preventing the accused from citing self-defence or even fear of attack as a justifiable reason for carrying a knife. The U.K. government advisory website on crime and justice formerly stated that "even if you carry a knife to protect yourself or make yourself feel safer but don’t intend to use it then you are committing a crime."
| Technology | Knives | null |
1137395 | https://en.wikipedia.org/wiki/Fishing%20net | Fishing net | A fishing net is a net used for fishing. Some fishing nets are also called fish traps, for example fyke nets. Fishing nets are usually meshes formed by knotting a relatively thin thread. Early nets were woven from grasses, flaxes and other fibrous plant material. Later cotton was used. Modern nets are usually made of artificial polyamides like nylon, although nets of organic polyamides such as wool or silk thread were common until recently and are still used.
History
Fishing nets have been used widely in the past, including by stone age societies. The oldest known fishing net is the net of Antrea, found with other fishing equipment in the Karelian town of Antrea, Finland, in 1913. The net was made from willow, and dates back to 8300 BC. Recently, fishing net sinkers from 27,000 BC were discovered in Korea, making them the oldest fishing implements discovered, to date, in the world. The remnants of another fishing net dates back to the late Mesolithic, and were found together with sinkers at the bottom of a former sea. Some of the oldest rock carvings at Alta (4200–500 BC) have mysterious images, including intricate patterns of horizontal and vertical lines sometimes explained as fishing nets. American Native Indians on the Columbia River wove seine nets from spruce root fibers or wild grass, again using stones as weights. For floats they used sticks made of cedar which moved in a way which frightened the fish and helped keep them together. With the help of large canoes, pre-European Maori deployed seine nets which could be over one thousand metres long. The nets were woven from green flax, with stone weights and light wood or gourd floats, and could require hundreds of men to haul.
Fishing nets are well documented in antiquity. They appear in Egyptian tomb paintings from 3000 BC. In ancient Roman literature, Ovid makes many references to fishing nets, including the use of cork floats and lead weights. Pictorial evidence of Roman fishing comes from mosaics which show nets. In a parody of fishing, a type of gladiator called retiarius was armed with a trident and a cast net. He would fight against a secutor or the murmillo, who carried a short sword and a helmet with the image of a fish on the front. Between 177 and 180 the Greek author Oppian wrote the Halieutica, a didactic poem about fishing. He described various means of fishing including the use of nets cast from boats, scoop nets held open by a hoop, and various traps "which work while their masters sleep". Here is Oppian's description of fishing with a "motionless" net:
The fishers set up very light nets of buoyant flax and wheel in a circle round about while they violently strike the surface of the sea with their oars and make a din with sweeping blow of poles. At the flashing of the swift oars and the noise the fish bound in terror and rush into the bosom of the net which stands at rest, thinking it to be a shelter: foolish fishes which, frightened by a noise, enter the gates of doom. Then the fishers on either side hasten with the ropes to draw the net ashore.
In Norse mythology the sea giantess Rán uses a fishing net to trap lost sailors. | Technology | Hunting and fishing | null |
1137424 | https://en.wikipedia.org/wiki/Loricariidae | Loricariidae | Loricariidae is the largest family of catfish (order Siluriformes), with over 90 genera and just over 680 species. Loricariids originate from freshwater habitats of Costa Rica, Panama, and tropical and subtropical South America. These fish are noted for the bony plates covering their bodies and their suckermouths. Several genera are sold as "plecos", notably the suckermouth catfish, Hypostomus plecostomus, and are popular as aquarium fish.
Common names
Members of the family Loricariidae are commonly referred to as loricariids, suckermouth catfishes, armoured catfish, or suckermouth armoured catfish. The name "plecostomus", and its shortened forms "pleco" and "plec", are used for many Loricariidae, since Plecostomus plecostomus (now called Hypostomus plecostomus) was one of the first loricariid species imported for the fish-keeping hobby.
Some loricariids are not normally considered "plecostomus", such as Farlowella catfish.
In their native range, these fish are known as cascudos or acarís.
L-numbers
Some types of loricariids are often referred to by their 'L-number'; this has become common since imports of loricariid catfish from South America often included specimens that had not been taxonomically described. Currently, L-numbers are used not only by fish-keeping enthusiasts, but also by biologists, since they represent a useful stopgap until a new species of fish is given a full taxonomic name. In some cases, two different L-numbered catfish have turned out to be different populations of the same species, while in other cases, multiple (but superficially similar) species have all been traded under a single L-number.
Taxonomy and evolution
Because of their highly specialized morphology, loricariids have been recognized as a monophyletic assemblage in even the earliest classifications of the Siluriformes, meaning they consist of a natural grouping with a common ancestor and all of its descendants. Loricariidae is one of seven families in the superfamily Loricarioidea, along with Amphiliidae, Trichomycteridae, Nematogenyidae, Callichthyidae, Scoloplacidae, and Astroblepidae. Some of these families also exhibit suckermouths or armor, although never in the same individual, as in loricariids.
This is the largest catfish family, including about 684 species in around 92 genera, with new species being described each year. However, this family is in flux, and revisions are likely. For example, the subfamily Ancistrinae is accepted as late as the 2006 edition of Nelson's Fishes of the World; it later becomes grouped as a tribe, because of its recognition as a sister group to the Pterygoplichthyini. Under Ambruster, six subfamilies are recognized: Delturinae, Hypoptopomatinae, Hypostominae, Lithogeneinae, Loricariinae, and Neoplecostominae.
Monophyly for the family is strongly supported, except, possibly, the inclusion of Lithogenes. Lithogenes is the only genus within the subfamily Lithogeneinae. This genus and subfamily, the most basal group in Loricariidae, is the sister group to the rest of the family. Neoplecostominae are the most basal group among the loricariids with the exception of Lithogeneinae. However, the genera of Neoplecostominae do not appear to form a monophyletic assemblage. The two subfamilies Loricariinae and Hypoptopomatinae appear to be generally regarded as monophyletic. However, the monophyly and composition of the other subfamilies are currently being examined and will likely be altered substantially in the future. The Hypostominae are the largest subfamily of Loricariidae. It is made up of five tribes. Four of the five tribes, Corymbophanini, Hypostomini, Pterygoplichthyini, and Rhinelepini, include about 24 genera. The fifth and largest tribe, Ancistrini (formerly recognized as its own subfamily), includes 30 genera.
Loricariid fossils are extremely rare. The earliest known definitive taxon is Taubateia from the Late Oligocene-Early Miocene in Brazil. The putative Cenomanian member Afrocascudo, initially described as the earliest loricariid catfish in 2024, might represent a juvenile obaichthyid lepisosteiform, possibly a junior synonym of Obaichthys, though this has been disputed based on the complete ossification of the bones indicating full maturity and the absence of important holostean characters. Within the superfamily Loricarioidea, the Loricariidae are the most derived; in this superfamily, the trend is toward increasingly complex jaw morphology, which may have allowed for the great diversification of the Loricariidae, which have the most advanced jaws.
Distribution and habitat
The family Loricariidae is vastly distributed over both the east and west sides of the Andes mountains, however, most species are generally restricted to small geographic ranges. They are primarily found in freshwater habitats of South America, but several loricariines and hypostomines are native to Panama, and two species (Fonchiiichthys uracanthus and Hemiancistrus aspidolepis) are native to Costa Rica. Species occur in swift-flowing streams from the lowlands up to in elevation. They can also be found in a variety of other freshwater environments. They can be found in torrential mountain rivers, quiet brackish estuaries, black acidic waters, and even in subterranean habitats.
Description and biology
This family has extremely variable color patterns and body shapes. Loricariids are characterized by bony plates covering their bodies, similar to the bony plates in callichthyids. (In Latin, lorica means corselet). These fish exhibit a ventral suckermouth, with papillae (small projections) on the lips. When present, the adipose fin usually has a spine at the forward edge. These fish have, when they are present, a unique pair of maxillary barbels. These fish have relatively long intestines due to their usually herbivorous or detrivorous diets. The body is characteristically flattened in this family. Taste buds cover almost the entire surface of the body and fin spines. Body lengths can range from in Nannoplecostomus eleonorae to over in Panaque, Acanthicus, and Pterygoplichthys.
One of the most obvious characteristics of the loricariids is the suckermouth. The modified mouth and lips allow the fish to feed, breathe, and attach to the substrate through suction. The lips were once believed to be unable to function as a sucker while respiration continued, as the inflowing water would cause the system to fail; however, respiration and suction can function simultaneously. Inflowing water passing under the sucker is limited to a thin stream immediately behind each maxillary barbel; the maxillae in loricariids support only small maxillary barbels and are primarily used to mediate the lateral lip tissue in which they are embedded, preventing failure of suction during inspiration. To achieve suction, the fish presses its lips against the substrate and expands its oral cavity, causing negative pressure.
Also, unlike most other catfishes, the premaxillae are highly mobile, and the lower jaws have evolved towards a medial position, with the teeth pointed rostroventrally; these are important evolutionary innovations. The fish rotates its lower and upper jaws to scrape the substrate. Of the two, the lower jaws are more mobile.
Loricariid catfishes have evolved several modifications of their digestive tracts that function as accessory respiratory organs or hydrostatic organs. These complex structures would have been independently evolved a number of times within the family. This includes an enlarged stomach in the Pterygoplichthyini, Hypostomus, and Lithoxus, a U-shaped diverticulum in Rhinelepini, and a ring-like diverticulum in Otocinclus. It may be noted that even loricariids with unmodified stomachs have a slight ability to breathe air.
Considerable sexual dimorphism occurs in this family, most pronounced during the breeding season. For example, in Loricariichthys, the male has a large expansion of its lower lip, which it uses to hold a clutch of eggs. Ancistrus males have snouts with fleshy tentacles. In loricariids, odontodes develop almost anywhere on the external surface of the body and first appear soon after hatching; odontodes appear in a variety of shapes and sizes and are often sexually dimorphic, being larger in breeding males. In most Ancistrini species, sharp evertible cheek spines (elongated odontodes) are often more developed in males and are used in intraspecific displays and combat.
Omega iris
Unusual for bony fish, many species have a modified iris called an omega iris. The dorsal segment of the iris expands downward over the pupil to form a loop which can expand and contract, called an iris operculum; when light levels are high, the pupil reduces in diameter and the loop expands to cover the center of the pupil, giving rise to a crescent-shaped, light-transmitting partial pupil. This feature gets its name from its similarity to an upside-down uppercase Greek letter omega (Ω). The origins of this structure are unknown, but breaking up the outline of the highly visible eye has been suggested to aid camouflage in what are often highly mottled animals. Species in the tribe Rhinelepini are an exceptional group among loricariids, having a circular iris. The presence or absence of the iris operculum can be used for identification of species in the subfamily Loricariinae.
Genetics
As of 2000, only 56 loricariid species have been cytogenically investigated. The basal diploid number of chromosomes is 2n = 54 in this family, but with a wide variation in the chromosome number in this fish group, ranging from 2n = 36 in the Loricariinae, Rineloricaria latirostris, to 2n = 96 in a species of Upsilodus (Hemipsilichthys). Most members of the Ancistrini and Pterygoplichthyini have 52 chromosomes. Karyotypic evolution by means of centric fusions and centric fissions seems to be a common feature among loricariids; this is demonstrated by a higher number of biarmed chromosomes in species with lower diploid number and many uniarmed chromosomes in species with higher diploid numbers. Studies conducted with representatives of some genera of Hypostominae showed, within this group, the diploid number ranges from 2n = 52 to 2n = 80. However, the supposed wide karyotypic diversity the family Loricariidae or the subfamily Hypostominae
would present is almost exclusively restricted to the genus Hypostomus, and the species from the other genera had a conserved diploid number. In some species, there is a ZZ/ZW sex-determination system.
Ecology
The suckermouth exhibited by these catfish allows them to adhere to objects in their habitats, even in fast-flowing waters. The mouth and teeth also are adapted to feed on a variety of foods, such as algae, invertebrates, and detritus. Some species, notably the Panaque, are known for xylophagy, or the ability to digest wood.
Most species of loricariids are nocturnal animals. Some species are territorial, while others, such as Otocinclus, prefer to live in groups.
Air-breathing is well known among many loricariids; this ability is dependent on the risk of hypoxia faced by a species; torrent-dwelling species tend to have no ability to breathe air, while low-land, pool-dwelling species, such as those of Hypostomus, have a great ability to breathe air. Pterygoplichthys is known for being kept out of water and sold alive in fish markets, surviving up to 30 hours out of water. Loricariids are facultative air breathers; they will only breathe air if under stress and will only use their gills in situations when oxygen levels are high. The dry season is a likely time for this; there would be little food in the stomach, which would allow its use for air breathing.
Loricariids exhibit a wide range of reproductive strategies, including cavity spawning, attachment of eggs on the underside of rocks, and egg-carrying. Parental care is usually good, and the male guards the eggs and sometimes the larvae. The eggs hatch after four to 20 days, depending on the species.
Three species known from subterranean habitats are true troglobites with reduced pigmentation (appearing overall whitish) and eyes: Ancistrus cryptophthalmus, A. galani and A. formoso. Similar adaptions with reduced pigmentation are known from two loricariids found in deep water in large Amazonian rivers, Peckoltia pankimpuju and Panaque bathyphilus.
In the aquarium
Loricariids are popular aquarium fish, where they are often sold as "plecs", "plecos" or "plecostomus". These fish are often purchased because of their algae-eating habits, though this role may not be carried out. Loricariid are either vegetarian, omnivore, carnivore or wood-eaters. A great many species of loricariids are also sold for their ornamental qualities, representing many body shapes and colors.
Most species of loricariids are nocturnal and will shy away from bright light, appreciating some sort of cover to hide under throughout the day. As they often originate from habitats with fast-moving water, filtration should be vigorous.
A number of species of loricariids have been bred in captivity.
| Biology and health sciences | Siluriformes | null |
4472066 | https://en.wikipedia.org/wiki/Atomic%20formula | Atomic formula | In mathematical logic, an atomic formula (also known as an atom or a prime formula) is a formula with no deeper propositional structure, that is, a formula that contains no logical connectives or equivalently a formula that has no strict subformulas. Atoms are thus the simplest well-formed formulas of the logic. Compound formulas are formed by combining the atomic formulas using the logical connectives.
The precise form of atomic formulas depends on the logic under consideration; for propositional logic, for example, a propositional variable is often more briefly referred to as an "atomic formula", but, more precisely, a propositional variable is not an atomic formula but a formal expression that denotes an atomic formula. For predicate logic, the atoms are predicate symbols together with their arguments, each argument being a term. In model theory, atomic formulas are merely strings of symbols with a given signature, which may or may not be satisfiable with respect to a given model.
Atomic formula in first-order logic
The well-formed terms and propositions of ordinary first-order logic have the following syntax:
Terms:
,
that is, a term is recursively defined to be a constant c (a named object from the domain of discourse), or a variable x (ranging over the objects in the domain of discourse), or an n-ary function f whose arguments are terms tk. Functions map tuples of objects to objects.
Propositions:
,
that is, a proposition is recursively defined to be an n-ary predicate P whose arguments are terms tk, or an expression composed of logical connectives (and, or) and quantifiers (for-all, there-exists) used with other propositions.
An atomic formula or atom is simply a predicate applied to a tuple of terms; that is, an atomic formula is a formula of the form P (t1 ,…, tn) for P a predicate, and the tn terms.
All other well-formed formulae are obtained by composing atoms with logical connectives and quantifiers.
For example, the formula ∀x. P (x) ∧ ∃y. Q (y, f (x)) ∨ ∃z. R (z) contains the atoms
.
As there are no quantifiers appearing in an atomic formula, all occurrences of variable symbols in an atomic formula are free.
| Mathematics | Mathematical logic | null |
7757921 | https://en.wikipedia.org/wiki/Brocade | Brocade | Brocade () is a class of richly decorative shuttle-woven fabrics, often made in coloured silks and sometimes with gold and silver threads. The name, related to the same root as the word "broccoli", comes from Italian meaning 'embossed cloth', originally past participle of the verb 'to stud, set with nails', from , 'small nail', from Latin , 'projecting, pointed'.
Brocade is typically woven on a draw loom. It is a supplementary weft technique; that is, the ornamental brocading is produced by a supplementary, non-structural, weft in addition to the standard weft that holds the warp threads together. The purpose of this is to give the appearance that the weave was actually embroidered on.
In Guatemala, brocade is the most popular technique used to decorate fabric woven by Maya weavers on backstrap looms.
Ornamental features in brocade are emphasised and wrought as additions to the main fabric, sometimes stiffening it, though more frequently producing on its face the effect of low relief. In some, but not all, brocades, these additions present a distinctive appearance on the back of the material where the supplementary weft or floating threads of the brocaded or broached parts hang in loose groups or are clipped away. When the weft is floating on the back, this is known as a continuous brocade; the supplementary weft runs from selvage to selvage. The yarns are cut away in cutwork and broché. Also, a discontinuous brocade is where the supplementary yarn is only woven in the patterned areas.
History
China
The manufacture of brocade began during the Warring States period of China. Many products of brocade have been found in tombs of the era. Several distinct styles of brocade have been developed in China, the most famous being Yunjin (Cloud brocade) of Nanjing, Song brocade of Suzhou, and Shu brocade of Chengdu.
Southeast Asia
Songket is a type of brocade in the Malay world (Indonesia, Malaysia, Singapore, Brunei etc.)
Byzantium
Dating back to the Middle Ages, brocade fabric was one of the few luxury fabrics worn by nobility throughout China, India, Persia, Greece, Japan, Korea and Byzantium. Woven by the Byzantines, brocades were an especially desirable fabric. From the 4th to the 6th centuries, production of silk was seemingly non-existent, as linen and wool were the predominant fabrics. During this period, there was no public knowledge of silk fabric production except for that which was kept secret by the Chinese. Over the years, knowledge of silk production became known among other cultures and spread westward. As silk production became known to Western cultures, trade from the East began to decrease. It was discovered by Byzantine historians that in the 6th century a pair of monks brought the secret of sericulture – silk production – to the Byzantine emperor. As a result, Western cultures were able to learn how to breed, raise, and feed silkworms. From this point until the 9th century, Byzantium became the biggest and most central producer for all of the Western world in the production of all types of silk motifs, including brocades, damasks, brocatelles and tapestry-like fabrics.
During the Early Middle Ages, brocade fabrics were available only to the wealthiest of people as the Byzantine emperor charged extreme prices for the fabric. The designs woven into brocade fabrics were often Persian in origin. It was also common to see Christian subjects depicted in the complex weaves of the fabric. When these luxurious fabrics were made into clothing or wall hangings, they were at times adorned with precious and semiprecious stones, small medallions of enamel, embroidery and appliqués.
The Late Middle Ages
Wealthy noblemen and noblewomen dressed in silk brocades from Italy, and velvets trimmed with fur from Germany. During the 14th and 15th centuries, the Court of Burgundy was made known for their continuous fashionable tastes and luxurious dress.
Renaissance Italy
Brocades were also an important fabric during the Renaissance, and especially the Italian Renaissance. As wool and silk were the primary fabrics used by Europeans during the Renaissance, and despite the lack of documentary evidence, it is said that due to the increase in complexity of decoration of Italian silk fabrics of the 15th century, there must have been improvements in silk-weaving looms around this time. The complexity and high quality of luxurious silk fabrics caused Italy to become the most important and superior manufacturer of the finest silk fabrics for all of Europe. The almost sculptural lines of the fashions during the Renaissance were paired perfectly with the exquisite beauty and elegance of brocade, damask, and other superior silk textiles.
Modern uses
Brocade fabrics are mostly for upholstery and draperies. They are also used for evening and formal clothing, for vestments, as well as for costumes. In India, Banarasi brocade is extensively used in decorating Banarasi saris, dresses, and dupattas. The use of precious and semi-precious stones in the adornment of brocades is not common, but has been replaced by the use of sequins and beading. Brocade fabrics are now largely woven on a Jacquard loom that is able to create many complex tapestry-like designs using the Jacquard technique. Although many brocade fabrics look like tapestries and are advertised in some fashion promotions as such, they are not to be confused with true tapestries. Patterns such as brocade, brocatelle, damask and tapestry-like fabrics are known as jacquard patterns.
Gallery
| Technology | Weaving | null |
7760747 | https://en.wikipedia.org/wiki/Theoretical%20plate | Theoretical plate | A theoretical plate in many separation processes is a hypothetical zone or stage in which two phases, such as the liquid and vapor phases of a substance, establish an equilibrium with each other. Such equilibrium stages may also be referred to as an equilibrium stage, ideal stage, or a theoretical tray. The performance of many separation processes depends on having series of equilibrium stages and is enhanced by providing more such stages. In other words, having more theoretical plates increases the efficiency of the separation process be it either a distillation, absorption, chromatographic, adsorption or similar process.
Applications
The concept of theoretical plates and trays or equilibrium stages is used in the design of many different types of separation.
Distillation columns
The concept of theoretical plates in designing distillation processes has been discussed in many reference texts. Any physical device that provides good contact between the vapor and liquid phases present in industrial-scale distillation columns or laboratory-scale glassware distillation columns constitutes a "plate" or "tray". Since an actual, physical plate can never be a 100% efficient equilibrium stage, the number of actual plates is more than the required theoretical plates.
where is the number of actual, physical plates or trays, is the number of theoretical plates or trays and is the plate or tray efficiency.
So-called bubble-cap or valve-cap trays are examples of the vapor and liquid contact devices used in industrial distillation columns. Another example of vapor and liquid contact devices are the spikes in laboratory Vigreux fractionating columns.
The trays or plates used in industrial distillation columns are fabricated of circular steel plates and usually installed inside the column at intervals of about 60 to 75 cm (24 to 30 inches) up the height of the column. That spacing is chosen primarily for ease of installation and ease of access for future repair or maintenance.
An example of a very simple tray is a perforated tray. The desired contacting between vapor and liquid occurs as the vapor, flowing upwards through the perforations, comes into contact with the liquid flowing downwards through the perforations. In current modern practice, as shown in the adjacent diagram, better contacting is achieved by installing bubble-caps or valve caps at each perforation to promote the formation of vapor bubbles flowing through a thin layer of liquid maintained by a weir on each tray.
To design a distillation unit or a similar chemical process, the number of theoretical trays or plates (that is, hypothetical equilibrium stages), , required in the process should be determined, taking into account a likely range of feedstock composition and the desired degree of separation of the components in the output fractions. In industrial continuous fractionating columns, is determined by starting at either the top or bottom of the column and calculating material balances, heat balances and equilibrium flash vaporizations for each of the succession of equilibrium stages until the desired end product composition is achieved. The calculation process requires the availability of a great deal of vapor–liquid equilibrium data for the components present in the distillation feed, and the calculation procedure is very complex.
In an industrial distillation column, the required to achieve a given separation also depends upon the amount of reflux used. Using more reflux decreases the number of plates required and using less reflux increases the number of plates required. Hence, the calculation of is usually repeated at various reflux rates. is then divided by the tray efficiency, E, to determine the actual number of trays or physical plates, , needed in the separating column. The final design choice of the number of trays to be installed in an industrial distillation column is then selected based upon an economic balance between the cost of additional trays and the cost of using a higher reflux rate.
There is a very important distinction between the theoretical plate terminology used in discussing conventional distillation trays and the theoretical plate terminology used in the discussions below of packed bed distillation or absorption or in chromatography or other applications. The theoretical plate in conventional distillation trays has no "height". It is simply a hypothetical equilibrium stage. However, the theoretical plate in packed beds, chromatography and other applications is defined as having a height.
The empirical formula known as Van Winkle's Correlation can be used to predict the Murphree plate efficiency for distillation columns separating binary systems.
Distillation and absorption packed beds
Distillation and absorption separation processes using packed beds for vapor and liquid contacting have an equivalent concept referred to as the plate height or the height equivalent to a theoretical plate (HETP). HETP arises from the same concept of equilibrium stages as does the theoretical plate and is numerically equal to the absorption bed length divided by the number of theoretical plates in the absorption bed (and in practice is measured in this way).
where is the number of theoretical plates (also called the "plate count"), is the total bed height and is the height equivalent to a theoretical plate.
The material in packed beds can either be random dumped packing (1-3" wide) such as Raschig rings or structured sheet metal. Liquids tend to wet the surface of the packing and the vapors contact the wetted surface, where mass transfer occurs.
Chromatographic processes
The theoretical plate concept was also adapted for chromatographic processes by Martin and Synge. The IUPAC's Gold Book provides a definition of the number of theoretical plates in a chromatography column.
The same equation applies in chromatography processes as for the packed bed processes, namely:
In packed column chromatography, the HETP may also be calculated with the Van Deemter equation.
In capillary column chromatography HETP is given by the Golay equation.
Other applications
The concept of theoretical plates or trays applies to other processes as well, such as capillary electrophoresis and some types of adsorption.
| Physical sciences | Separation processes | Chemistry |
88580 | https://en.wikipedia.org/wiki/Boidae | Boidae | The Boidae, commonly known as boas or boids, are a family of nonvenomous snakes primarily found in the Americas, as well as Africa, Europe, Asia, and some Pacific islands. Boas include some of the world's largest snakes, with the green anaconda of South America being the heaviest and second-longest snake known; in general, adults are medium to large in size, with females usually larger than the males. Six subfamilies comprising 15 genera and 54 species are currently recognized.
Description
Like the pythons, boas have elongated supratemporal bones. The quadrate bones are also elongated, but not as much, while both are capable of moving freely so when they swing sideways to their maximum extent, the distance between the hinges of the lower jaw is greatly increased.
Both families share a number of primitive characteristics. Nearly all have a relatively rigid lower jaw with a coronoid element, as well as a vestigial pelvic girdle with hind limbs that are partially visible as a pair of spurs, one on either side of the vent. In males, these anal spurs are larger and more conspicuous than in females. A long row of palatal teeth is present, and most species have a functional left lung that can be up to 75% as large as the right lung.
Boids are, however, distinguished from the pythons in that none has postfrontal bones or premaxillary teeth, and that they give birth to live young. When labial pits are present, these are located between the scales as opposed to on them. Also, their geographical distributions are almost entirely mutually exclusive. In the few areas where they do coexist, the tendency is for them to occupy different habitats.
Formerly, boas were said to be found in the New World and pythons in the Old World. While this is true of boine boas, other boid species are present in Africa, much of southern Eurasia, Madagascar, New Guinea, and the Solomon Islands, so this is not accurate. However, they seem more abundant in evolutionarily isolated areas. South America was isolated until a few million years ago, with a fauna that included marsupials and other distinctive mammals. With the formation of the Panamanian land bridge to North America about three million years ago, boines have migrated north as colubrids (and various Nearctic mammals) have migrated south, as part of the Great American Interchange.
Distribution and habitat
Most species are found in North, Central, and South America, as well as the Caribbean, while a few are found in southeastern Europe and Asia Minor, North, Central and East Africa, Madagascar, the Arabian Peninsula, Central and Southwestern Asia, India and Sri Lanka, Indonesian islands (Moluccas, West Papua, Talaud, Sulawesi) and Papua New Guinea through Melanesia and Samoa.
Feeding
Prey is killed by constriction; after an animal has been grasped to restrain it, a number of coils are hastily wrapped around it. Then, by applying and maintaining sufficient pressure, the snake prevents its prey from inhaling, so that it eventually succumbs to asphyxiation. Recently, the pressures produced during constriction have been suggested as the cause of cardiac arrest by interfering with blood flow, but this hypothesis has not yet been confirmed.
Larger specimens usually eat animals about the size of a domestic cat, but larger food items are not unknown: the diet of the green anaconda (Eunectes murinus) is known to include subadult tapirs. Prey is swallowed whole, and may take several days or even weeks to fully digest. Despite their intimidating size and muscular power, they are generally not dangerous to humans.
Contrary to popular belief, even the larger species do not crush their prey to death; in fact, prey is not even noticeably deformed before it is swallowed. The speed with which the coils are applied is impressive and the force they exert may be significant, but death is caused by suffocation, with the victim not being able to move its ribs to breathe while it is being constricted.
Reproduction
Most species of boa are ovoviviparous, with females giving birth to live young. This is in contrast to the pythons, which lay eggs (oviparous).
Subfamilies
Type genus = Boa – Gray, 1825
Taxonomy
Pythons were historically classified as a subfamily of Boidae (called Pythoninae), but it was later determined that they are not closely related to boas despite having superficial similarities.
Almost all of the non-boine boids are frequently elevated to their own full families: Calabariidae/inae, Candoiidae/inae, Charinidae/inae, Erycidae/inae, Sanziniidae/inae, and Ungaliophiidae/inae. The taxonomy of boid snakes has been long debated, and ultimately the decision whether to assign a particular clade to a particular Linnaean rank (such as a superfamily, family, or subfamily) is arbitrary.
The subfamily Ungaliophiinae was formerly made up of four genera. Two of them (Tropidophis and Trachyboa) are actually more closely related to the American pipe snake (Anilius scytale) than to the boas, and are now placed in the family Tropidophiidae within the superfamily Amerophidia. The other two genera (Ungaliophis and Exiliboa) are the sister group of the Charina/Lichanura clade within Boidae.
Gallery
| Biology and health sciences | Reptiles | null |
88590 | https://en.wikipedia.org/wiki/Megatherium | Megatherium | Megatherium ( ; from Greek () 'great' + () 'beast') is an extinct genus of ground sloths endemic to South America that lived from the Early Pliocene through the end of the Late Pleistocene. It is best known for the elephant-sized type species Megatherium americanum, primarily known from the Pampas, but ranging southwards to northernmost Patagonia and northwards to southern Bolivia during the late Middle Pleistocene and Late Pleistocene. Various other species belonging to the subgenus Pseudomegatherium ranging in size comparable to considerably smaller than M. americanum are known from the Andean region.
The first (holotype) specimen of Megatherium was discovered in 1787 on the bank of the Luján River in what is now northern Argentina. The specimen was then shipped to Spain the following year wherein it caught the attention of the French paleontologist Georges Cuvier, who named the animal in 1796 and was the first to determine, by means of comparative anatomy, that Megatherium was a giant sloth.
Megatherium is part of the sloth family Megatheriidae, which also includes the closely related and similarly giant Eremotherium, comparable in size to M. americanum, which was native to tropical South America, Central America and North America as far north as the southern United States.
Megatherium americanum is thought to have been a browser that fed on the foliage and twigs of trees and shrubs using a black rhinoceros–like prehensile upper lip. Despite its large body size, Megatherium americanum is widely thought to have been able to adopt a bipedal posture at least while standing, which allowed it to feed on high-growing leaves, as well as possibly to use its claws for defense.
Megatherium became extinct around 12,000 years ago as part of the end-Pleistocene extinction event, simultaneously with the majority of other large mammals in the Americas. The extinctions followed the first arrival of humans in the Americas, and one and potentially multiple kill sites where M. americanum was slaughtered and butchered is known, suggesting that hunting could have been a factor in its extinction.
Research history
The earliest specimen of Megatherium americanum was discovered in 1787 by Manuel de Torres, a Dominican friar and naturalist, from a ravine on the banks of the Lujan River in what is now northern Argentina, which at the time was part of the Viceroyalty of the Río de la Plata in the Spanish Empire. Torres described the bones as a ‘wonder and providence of the Lord’. On the orders of the then viceroy of la Plata, Nicolás Cristóbal del Campo, Marqués de Loreto, the specimen was moved to the capital Buenos Aires. There the skeleton was drawn for the first time by José Custodio Sáa y Faria in a horse-like posture. Campo summoned a number of local indigenous leaders to ask if they had heard of the animal. The skeleton was then transferred by Campo to the Royal Cabinet of Natural History of Madrid (now the National Museum of Natural Sciences MNCN) in 7 crates, which had arrived and been unpacked by late 1788.
At the direction of the cabinets main taxidermist Juan Bautista Bru, the specimen was then mounted for public exhibition (which remains unaltered in the modern museum display). In 1796 a scientific description of the skeleton was published authored by Bru along with engineer Joseph Garriga, with engravings by Manuel Navarro. As the work was going through the process of publication in 1795, preliminary prints of the paper were obtained by French diplomat Philippe-Rose Roume who was in Madrid at the time, who sent them to the National Museum of Natural History (Muséum national d'histoire naturelle) in Paris, France, where they were seen by French anatomist and paleontologist Georges Cuvier.
Cuvier, working solely from the prints from Madrid and not visiting the specimen personally, and using comparative anatomy with "edentate" mammals (now recognised as members of the order Xenarthra) in the collection of the Paris museum, correctly recognised that the remains represented those of a giant sloth, and an animal that was entirely extinct and not living. In early 1796, somewhat before the full publication of the work by Bru, Garriga and Navarro, Cuvier published a paper naming the species Megatherium americanum (literally "Great American beast"), becoming the first fossil mammal to be identified with both a genus and species name. Which description had priority has been controversial in the past. Cuvier later wrote a fuller description in 1804, which was republished in his famous 1812 book Recherches sur les ossemens fossiles de quadrupèdes. Cuvier identified Megatherium as a sloth primarily on the basis of its skull morphology, the dental formula and the shoulder, while regarding the anatomy of its limbs as more similar to armadillos and anteaters. Cuvier suggested that based on the proportions of its limbs (which are approximately equal to each other), that Megatherium did not jump or run, nor crawl like living sloths, with the presence of a clavicle and well developed crests on the humerus, suggesting to Cuvier that the animal probably used its forelimbs to grasp. A later publication in 1823 by Cuvier suggested that giant carapaces found in the Pampas also belonged to Megatherium, but British paleontologist Richard Owen in 1839 demonstrated that these actually belonged to another extinct group of xenarthrans called glyptodonts that were related to armadillos.
Additional remains of Megatherium were collected by Charles Darwin during the Voyage of the Beagle in the 1830s, these remains were assigned by Richard Owen in 1840 to the species Megatherium cuvieri, which had been named by Anselme Gaëtan Desmarest in 1822. These remains are now assigned to M. americanum.Owen later wrote a monograph series from 1851 to 1860 thoroughly describing the anatomy of M. americanum.
From the late 19th century onward additional species of Megatherium were described. In 1888 Argentine explorer Francisco Moreno erected the species Megatherium filholi for remains found in the Late Pleistocene of Argentina. In 1880 Paul Gervais and Florentino Ameghino described the species M. tarijense from remains of Pleistocene age found in Bolivia. In 1893 Rodolfo Amando Philippi erected the species M. sundti and M. medinae from remains found in the Pleistocene of Bolivia and Chile, respectively. In 1921, Florentino's brother Carlos Ameghino and Lucas Kraglievich described the species Megatherium gallardoi based on remains found in the Pampas of Northern Argentina, of Early-Middle Pleistocene age. In 2001, the species M. altiplanicum was described based on remains found in the Pliocene of Bolivia. In 2004, the species Megatherium urbinai was erected based on remains found in Pleistocene aged deposits in Peru. In 2006, the species Megatherium celendinense was erected for remains of Pleistocene age found in the Peruvian Andes.
Taxonomy and evolution
Megatherium is divided into 2 subgenera, Megatherium and Pseudomegatherium. Taxonomy according to Pujos (2006) and De Iuliis et al (2009):
Subgenus Megatherium
†M. altiplanicum
†M. americanum
†M. gallardoi Ameghino & Kraglievich, 1921
Subgenus Pseudomegatherium
†M. celendinense
†M. medinae
†M. sundti
†M. tarijense
†M. urbinai
Megatherium gallardoi Ameghino & Kraglievich, 1921 from the Pampas dating to the Early to Middle Pleistocene has sometimes been regarded as a synonym of M. americanum. The species Megatherium filholi Moreno, 1888 also from the Pleistocene of the Pampas region, historically regarded to be a junior synonym of M. americanum representing juvenile individuals has been suggested to be valid by some recent authors. Megatherium gaudryi Moreno (1888) from Argentina, of uncertain temporal provenance but possibly Pliocene in age, may also be valid.
Mitochondrial DNA sequences obtained from M. americanum indicates that three-toed sloths (Bradypus) are their closest living relatives. Phylogeny of sloths after Delsuc et al. 2019.
Megatheriidae is suggested to have diverged from other sloth families during the Oligocene, around 30 million years ago. The subfamily to which Megatherium belongs, Megatheriinae, first appeared in the Middle Miocene in Patagonia, at least 12 million years ago, represented by the genus Megathericulus. The earliest known remains of the genus Megatherium are known from the Pliocene, found in Bolivia (M. altiplanicum) and the Pampas (indeterminate species), dating to at least 3.6 million years ago. M. altiplanicum is suggested to be more closely related to M. americanum than to species of Pseudomegatherium. Phylogeny of Megatheriinae after Pujos, 2006:
Megatherium americanum first appears in the fossil record during the second half of the Middle Pleistocene, from around 400,000 years ago.
Description
Size
M. americanum is one of the largest known ground sloths, with a total body length of around . Volumetric analysis suggests that a full grown M. americanum weighed around , comparable to an Asian elephant. The Late Pleistocene Andean-Altiplano Pseudomegatherium species Megatherium celendinense was likely comparable in size. These species were only rivalled in size amongst ground sloths by the closely related Eremotherium and the distantly related Lestodon. The Chilean Pseudomegatherium species M. sundti was much smaller, with an estimated body mass of only , with the Peruvian Megatherium urbinai, Bolivian Megatherium tarijense and the Chilean Megatherium medinae (all also belonging to Pseudomegatherium) also having a considerably smaller body size than M. americanum. The Pliocene Megatherium (Megatherium) species M. altiplanicum has been estimated to weigh .
Skull and jaws
The head of Megatherium is relatively small compared to body size. The skull of M. americanum has a relatively narrow snout/muzzle with a ossified nasal septum, and is suggested to have had a thick prehensile upper lip, similar to that of the living black rhinoceros. The morphology of the hyoid bones in Megatherium suggests that they were relatively rigid, this along with the short distance between the hyoid and the mandibular symphysis (the joint connecting the two halves of the lower jaw) suggests that the tongue had limited ability to protrude, and thus Megatherium lacked the long prehensile tongue often attributed to it historically. The skull is roughly cylindrical in shape, with the cranial region of the skull being narrow. The jugal bone of M. americanum has strongly developed ascending and descending processes. The skull of M. americanum has a relatively small cranial cavity (and thus brain) relative to skull size, with the skull having extensive sinus spaces. In many species of Megatherium, the lower jaw is relatively deep, which served to accommodate the very long hypselodont (evergrowing) teeth, which are considerably proportionally longer than those of other ground sloths. Like other ground sloths, the number of teeth in the jaw is reduced to 5 and 4 teeth in each half of the upper and lower jaws, respectively, and the teeth lack enamel. The teeth of Megatherium americanum have sharp crests separated by v-shaped valleys, which interlock with the teeth on the opposing jaw. These teeth were self-sharpening, akin to rodent incisors. The skull and jaws of M. americanum show adaptation to powerful vertical biting. M. americanum and M. altiplanicum are distinguished from species of the subgenus Pseudomegatherium by the fusion of the maxilla and premaxilla, while members of Pseudomegatherium are distinguished from those species by their flat occipital condyles.
Axial skeleton
Like other xenarthrans, the posterior trunk vertebrae of Megatherium americanum have additional xenarthrous processes that articulate with the other vertebrae. The ischium was connected to the caudal vertebrae, forming a synsacrum. The sacrum was composed of 5 vertebrae. The pubic symphysis is reduced. The tail is large in size.
Limbs
The bones of the forelimbs of M. americanum are relatively elongate and thin. The three fingers in the middle of the hand bore claws, while the cuneiform hand bones did not touch the ulna. The olecranon process of the ulna was relatively short. Like other xenarthrans, but unlike most other mammals, Megatherium possesses clavicles (collarbones), which serves to support the forelimb. Like other sloths, the clavicle is merged with the acromion of the scapula. The femur was massive and roughly rectangular in shape. As in most megatheriines, the tibia and fibula of Megatherium species are fused together at their proximal (closest to hip) end, while in M. americanum and M. tarijense, they are also fused together at their distal (closest to foot) ends. The foot was heavily modified from those of other mammals and earlier ground sloths, with a reduction in the number of digits on the inner part of the foot (digits I and II being lost), the increase in the size and robustness (thickness) of the metapodial elements of the outer digits, with the loss or reduction of the phalangeal bones. The calcaneum is wide and elongate posteriorly. The foot is suggested to have been inwardly rotated, historically the foot was suggested to be near vertical, though a recent study suggests that the angle was much shallower. The weight was primarily borne on the outer digits and the calcaneum. M. urbinai differs from M. americanum based on various characters of the feet and hands.
Ecology
Remains of Megatherium americanum have been found in low elevation areas to the east of the Andes mountains in northern Patagonia, the Pampas and adjacent areas in what is now northern Argentina, Uruguay, Paraguay, southern Bolivia and Rio Grande do Sul in southern Brazil. Megatherium americanum inhabited temperate, arid-to semi arid open habitats. During the Last Glacial Period, the Pampas was generally drier than it is at present with many areas exhibiting a steppe-like environment dominated by grass, with some areas of woodland.
Although some authors have suggested that Megatherium was an omnivore, isotopic analysis has supported an entirely herbivorous diet for Megatherium. Megatherium americanum is suggested to have been a browser that was a selective feeder on the foliage, twigs and fruits of trees and shrubs. The sharp cusps of the teeth served to shear plant material. Megatherium is widely thought to have been able to adopt a bipedal posture to use its forelimbs to grasp vegetation, though whether it was capable of moving in this posture is uncertain. Analysis of injuries on the clavicles of M. americanum individuals suggests that the species probably habitually moved in a quadrupedal posture and assumed a bipedal posture next to trees to feed on high-growing leaves, likely using its forelimbs to brace itself against the tree trunk, as well as to pull down higher branches within reach of its prehensile lip. Isotopic analysis suggests that some individuals of M. americanum at certain times and places also consumed grass. The smaller Megatherium tarijense has been suggested to have had a mixed feeding-browsing diet. Preserved coprolites attributed to Megatherium suggests that its diet included plants like Fabiana, Ephedra (Ephedra breana), beebrush, Junellia, and Chuquiraga.
Whether or not Megatherium had a slow metabolism like living tree sloths is uncertain. Analysis of the nutrient foramina in the diaphysis (shaft) of the femur of Megatherium americanum shows that they are more similar to those of other large living mammals like elephants than living tree sloths, which may suggest that it had a metabolism more similar to non-xenarthran mammals and was capable of vigorous activity similar to living elephants. However, isotopic analysis of teeth suggests that Megatherium had a somewhat lower body temperature than non-xenarthran mammals, around , comparable to that of living tree sloths, implying a lower metabolic rate. Megatherium americanum has been traditionally reconstructed as being covered with a thick coat of fur. Due to its very large body size, some authors have alternatively argued that Megatherium americanum was probably hairless like modern elephants for thermodynamic reasons. However this has been disputed, with other authors suggesting based on thermodynamic modelling assuming a living xenarthran-like metabolism that Megatherium species probably had a dense coat of fur around thick to be able to tolerate the relatively cool environments they inhabited.
Based on fossil trackways and the anatomy of its inner ear, which is considerably different from living sloths and more similar to those of armadillos, species of Megatherium, while probably not capable of moving at considerable speed due to limitations of their skeletal anatomy (with one study estimating a max speed of approximately or , a fraction of the or top speed observed for living elephants) were likely significantly more agile and mobile than living sloths, which are only capable of moving . Species of Megatherium likely relied on their large adult body size to protect themselves against predators. Like many other large mammals, Megatherium is suggested to have had a slow life cycle in accordance with a K-selection strategy. Megatherium americanum is suggested to have given birth to a single large offspring at a time.
The anatomy of its forelimb bones suggests that M. americanum had the ability to rapidly and powerfully extend its arms, which likely made its claws effective stabbing weapons. It may have used its claws like this to defend itself, as living tree sloths do. Although some authors in the 19th century suggested that Megatherium engaged in digging behaviour, this has been disputed by other scholars, and the morphology of its limb bones do not appear to display significant adaptations to digging unlike some ground sloths such as mylodontids.
In the Pampas, Megatherium americanum lived alongside other megafauna species, including the large ground sloth Lestodon, along with the smaller (but still large) ground sloths Mylodon, Glossotherium, and Scelidotherium, the glyptodonts (very large armadillos with fused round carapaces covering the body) Glyptodon, Doedicurus, and Panochthus, the large camel-like ungulate Macrauchenia and rhinoceros-like Toxodon, the gomphothere (elephant-relative) Notiomastodon, the equines Hippidion and Equus neogeus, the large short-faced bear Arctotherium, and the large sabertooth cat Smilodon. The range of Megatherium americanum overlaps little with its similarly sized tropical relative Eremotherium, with their co-occurrence only confidently reported from a few localities in Southern Brazil, and it is unclear whether they were contemporary at these localities.
Relationship with humans and extinction
During the Late Pleistocene, six species of Megatherium were present in South America, including M. americanum in the Pampas and adjacent regions, and the 5 species of Pseudomegatherium in the vicinity of the Andes.
The youngest unambiguous dates for Megatherium are from the end of the Late Pleistocene. Supposed early Holocene dates obtained for Megatherium americanum and other Pampas megafauna have been questioned, with suggestions that they are likely due to humic acid contamination of the collagen used to radiocarbon date the bones. Megatherium disappeared simultaneously along with the vast majoriy (>80%) of other large (megafaunal) South American mammals, as part of the end-Pleistocene extinction event. The use of bioclimatic envelope modeling indicates that the area of suitable habitat for Megatherium had shrunk and become fragmented by the mid-Holocene. While this alone would not likely have caused its extinction, it has been cited as a possible contributing factor.
Towards the end of the Late Pleistocene, humans first arrived in the Americas, with some of the earliest evidence of humans in South America being the Monte Verde II site in Chile, dating to around 14,500 years Before Present (~12,500 BC). The extinction interval of Megatherium and other megafauna coincides with the appearance and abundance of Fishtail points, which are suggested to have been used to hunt megafauna, across the Pampas region and South America more broadly. At the Paso Otero 5 site in the Pampas of northeast Argentina, Fishtail points are associated with burned bones of Megatherium americanum and other extinct megafauna. The bones appear to have been deliberately burned as a source of fuel. Due to the poor preservation of the bones there is no clear evidence of human modification.
There is evidence for the butchery of Megatherium by humans. Two M. americanum bones, an ulna and an atlas vertebra, from separate collections, bear cut marks suggestive of butchery, with the latter suggested to represent an attempt to exploit the contents of the head. A kill site dating to around 12,600 years Before Present (BP), is known from Campo Laborde in the Pampas in Argentina, where a single individual of M. americanum was slaughtered and butchered at the edge of a swamp, which is the only confirmed giant ground-sloth kill site in the Americas. At the site several stone tools were present, including the fragment of a projectile point. Another possible kill site is Arroyo Seco 2 near Tres Arroyos in the Pampas in Argentina, where M. americanum bones amongst those of other megafauna were found associated with human artifacts dating to approximately 14,782–11,142 cal yr BP. This hunting may have been a factor in its extinction.
Cultural references
The Megatherium Club, named for the extinct animal and founded by William Stimpson, was a group of Washington, D.C.–based scientists who were attracted to that city by the Smithsonian Institution's rapidly growing collection, from 1857 to 1866.
| Biology and health sciences | Xenarthra | Animals |
88595 | https://en.wikipedia.org/wiki/Reticulated%20python | Reticulated python | The reticulated python (Malayopython reticulatus) is a python species native to South and Southeast Asia. It is the world's longest snake, and the third heaviest after the green anaconda and Burmese python. It is listed as least concern on the IUCN Red List because of its wide distribution. In several countries in its range, it is hunted for its skin, for use in traditional medicine, and for sale as pets. Due to this, reticulated pythons are one of the most economically important reptiles worldwide.
It is an excellent swimmer, has been reported far out at sea, and has colonized many small islands within its range.
Like all pythons, it is a non-venomous constrictor. In very rare cases, adult humans have been killed (and in at least six reported cases, eaten) by reticulated pythons.
Taxonomy
The reticulated python was first described in 1801 by German naturalist Johann Gottlob Theaenus Schneider, who described two zoological specimens held by the Göttingen Museum in 1801 that differed slightly in colour and pattern as separate species—Boa reticulata and Boa rhombeata. The specific name, reticulatus, is Latin meaning "net-like", or reticulated, and is a reference to the complex color pattern. The generic name Python was proposed by French naturalist François Marie Daudin in 1803. American zoologist Arnold G. Kluge performed a cladistics analysis on morphological characters and recovered the reticulated python lineage as sister to the genus Python, hence not requiring a new generic name in 1993.
In a 2004 genetics study using cytochrome b DNA, Robin Lawson and colleagues discovered the reticulated python as sister to Australo-Papuan pythons, rather than Python molurus and relatives. Raymond Hoser erected the genus Broghammerus for the reticulated python in 2004, naming it after German snake expert Stefan Broghammer, on the basis of dorsal patterns distinct from those of the genus Python, and a dark mid-dorsal line from the rear to the front of the head, and red or orange (rather than brown) iris colour. In 2008, Lesley H. Rawlings and colleagues reanalysed Kluge's morphological data and combined it with genetic material, finding the reticulated clade to be an offshoot of the Australo-Papuan lineage as well. They adopted and redefined the genus name Broghammerus.
Most taxonomists choose to ignore Broghammerus and other names by Hoser as its description lacked scientific rigour and was not published in a reputable journal. R. Graham Reynolds and colleagues accordingly proposed the name Malayopython for this species and its sister species, the Timor python. Malayopython has been recognized by subsequent authors and the Reptile Database. Hoser has argued that Broghammerus was validly published and Malayopython name is invalid as it is a junior synonym. In 2021, the International Commission on Zoological Nomenclature found no basis for regarding the name Broghammerus to be invalid. Nevertheless, the name Malayopython remains in use by reliable sources.
Subspecies
Three subspecies have been proposed:
M. r. reticulatus (Schneider, 1801) – Asian reticulated python
M. r. jampeanus Auliya et al., 2002 – Kayaudi reticulated python or Tanahjampean reticulated python, about half the length, or according to Auliya et al. (2002), not reaching much more than in length. Found on Tanahjampea in the Selayar Archipelago south of Sulawesi. Closely related to M. r. reticulatus of the Lesser Sundas.
M. r. saputrai Auliya et al., 2002 – Selayer reticulated python, occurs on Selayar Island in the Selayar Archipelago and also in adjacent Sulawesi. This subspecies represents a sister lineage to all other populations of reticulated pythons tested. According to Auliya et al. (2002) it does not exceed in length.
The latter two are dwarf subspecies. Apparently, the population of the Sangihe Islands north of Sulawesi represents another such subspecies, which is basal to the P. r. reticulatus plus P. r. jampeanus clade, but it is not yet formally described.
The proposed subspecies M. r. "dalegibbonsi", M. r. "euanedwardsi", M. r. "haydnmacphiei", M. r. "neilsonnemani", M. r. "patrickcouperi", and M. r. "stuartbigmorei" have not found general acceptance.
Characteristics
The reticulated python has smooth dorsal scales that are arranged in 69–79 rows at midbody. Deep pits occur on four anterior upper labials, on two or three anterior lower labials, and on five or six posterior lower labials.
The reticulated python is the largest snake native to Asia. More than a thousand wild reticulated pythons in southern Sumatra were studied, and estimated to have a length range of , and a weight range of . Reticulated pythons with lengths more than are rare, though according to the Guinness Book of World Records, it is the only extant snake to regularly exceed that length. One of the largest scientifically measured specimens, from Balikpapan, East Kalimantan, Indonesia, was measured under anesthesia at and weighed after not having eaten for nearly 3 months.
The specimen once widely accepted as the largest-ever "accurately" measured snake, that being Colossus, a specimen kept at the Highland Park Zoo (now the Pittsburgh Zoo and Aquarium) in Pittsburgh, Pennsylvania, during the 1950s and early 1960s, with a peak reported length of from a measurement in November 1956, was later shown to have been substantially shorter than previously reported. When Colossus died on 14 April 1964, its body was deposited in the Carnegie Museum of Natural History. At that time, its skeleton was measured and found to be in total length, and the length of its fresh hide was measured as . The hide tends to stretch from the skinning process, thus may be longer than the snake from which it came – e.g., by roughly 20–40% or more. The previous reports had been constructed by combining partial measurements with estimations to compensate for "kinks", since completely straightening an extremely large live python is virtually impossible. Because of these issues, a 2012 journal article concluded, "Colossus was neither the longest snake nor the heaviest snake ever maintained in captivity." Too large to be preserved with formaldehyde and then stored in alcohol, the specimen was instead prepared as a disarticulated skeleton. The hide was sent to a laboratory to be tanned, but it was either lost or destroyed, and now only the skull and selected vertebrae and ribs remain in the museum's collection. Considerable confusion exists in the literature over whether Colossus was male or female (females tend to be larger).
Numerous reports have been made of larger snakes, but since none of these was measured by a scientist nor any of the specimens deposited at a museum, they must be regarded as unproven and possibly erroneous. In spite of what has been, for many years, a standing offer of a large financial reward (initially $1,000, later raised to $5,000, then $15,000 in 1978 and $50,000 in 1980) for a live, healthy snake or longer by the New York Zoological Society (later renamed as the Wildlife Conservation Society), no attempt to claim this reward has ever been made.
Reported sizes
The colour pattern is a complex geometric pattern that incorporates different colours. The back typically has a series of irregular diamond shapes flanked by smaller markings with light centers. In this species' wide geographic range, much variation of size, colour, and markings commonly occurs.
In zoo exhibits, the colour pattern may seem garish, but in a shadowy jungle environment amid fallen leaves and debris, it allows them to virtually disappear. Called disruptive colouration, it protects them from predators and helps them to catch their prey.
The huge size and attractive pattern of this snake has made it a favorite zoo exhibit, with several individuals claimed to be above in length and more than one claimed to be the largest in captivity. However, due to its huge size, immense strength, aggressive disposition, and the mobility of the skin relative to the body, it is very difficult to get exact length measurements of a living reticulated python, and weights are rarely indicative, as captive pythons are often obese. Claims made by zoos and animal parks are sometimes exaggerated, such as the claimed snake in Indonesia which was subsequently proven to be about long. For this reason, scientists do not accept the validity of length measurements unless performed on a dead or anesthetized snake that is later preserved in a museum collection or stored for scientific research.
A reticulated python kept in the United States in Kansas City, Missouri, named "Medusa" is considered by the Guinness Book of World Records to be the longest living snake ever kept in captivity. In 2011 it was reported to measure and weigh .
In 2012, an albino reticulated python, named "Twinkie", housed in Fountain Valley, California, was considered to be the largest albino snake in captivity by the Guinness World Records. It measured in length and weighed about .
Dwarf forms of reticulated pythons also occur, from some islands northwest of Australia, and these are being selectively bred in captivity to be much smaller, resulting in animals often referred to as "super dwarfs". Adult super dwarf reticulated pythons are typically between in length.
Distribution and habitat
The reticulated python is found in South and Southeast Asia from the Nicobar Islands, India, Bangladesh, Myanmar, Thailand, Laos, Cambodia, Vietnam, Malaysia, and Singapore, east through Indonesia and the Indo-Australian Archipelago (Sumatra, the Mentawai Islands, the Natuna Islands, Borneo, Sulawesi, Java, Lombok, Sumbawa, Sumba, Flores, Timor, Maluku, Tanimbar Islands) and the Philippines (Basilan, Bohol, Cebu, Leyte, Luzon, Mindanao, Mindoro, Negros, Palawan, Panay, Polillo, Samar, Tawi-Tawi). The original description does not include a type locality. The type locality was restricted to "Java" by Brongersma (1972).
Three subspecies have been proposed, but are not recognized in the Integrated Taxonomic Information System. The color and size can vary a great deal among the subspecies described. Geographical location is a good key to establishing the subspecies, as each one has a distinct geographical range.
The reticulated python lives in rainforests, woodlands, and nearby grasslands. It is also associated with rivers and is found in areas with nearby streams and lakes. An excellent swimmer, it has even been reported far out at sea and has consequently colonized many small islands within its range. During the early years of the 20th century, it is said to have been common even in busy parts of Bangkok, sometimes eating domestic animals.
Behaviour and ecology
Diet
As with all pythons, the reticulated python is an ambush predator, usually waiting until prey wanders within strike range before seizing it in its coils and killing by constriction. Its natural diet includes mammals and occasionally birds. Small specimens up to long eat mainly small mammals such as rats, other rodents, mouse-eared bats, and treeshrews, whereas larger individuals switch to prey such as small Indian civet and binturong, primates, pigs, and deer weighing more than . As a rule, the reticulated python seems able to swallow prey up to one-quarter its own length and up to its own weight. Near human habitation, it is known to snatch stray chickens, cats, and dogs on occasion.
Among the largest documented prey items are a half-starved sun bear of that was eaten by a specimen and took some 10 weeks to digest.
At least one case is reported of a foraging python entering a forest hut and taking a child.
Reproduction
The reticulated python is oviparous. Adult females lay between 15 and 80 eggs per clutch. At an optimum incubation temperature of , the eggs take an average of 88 days to hatch. Hatchlings are at least in length.
Danger to humans
The reticulated python is among the few snakes that prey on humans, and is the only species of snake where video and photographic proof exists of them having consumed humans. In 2015, the species was added to the Lacey Act of 1900, prohibiting import and interstate transport due to its "injurious" history with humans. Attacks on humans in captivity are not common. Wild pythons, however, are known to sometimes prey on humans, particularly in their natural habitat in Sulawesi, Indonesia. Considering the known maximum prey size, a full-grown reticulated python can open its jaws wide enough to swallow a human, but the width of the shoulders of some adult Homo sapiens can pose a problem for even a snake with sufficient size. Reports of human fatalities and human consumption (the latest examples of consumption of an adult human being well authenticated) include:
A report of a visit of Antonio van Diemen, Governor-General of the Dutch East India Company, to the Banda Islands in 1638, includes a description of an enslaved woman who, when tending to a garden on the volcanic island of Gunung Api, was strangled by a snake of "24 houtvoeten" (slightly over seven meters) in length, and then swallowed whole. The snake, having become slow after ingesting such a large prey, was subsequently shot by Dutch soldiers and brought to the Governor-General to be looked at, with its victim still inside. Although less reliable than this first-hand document, several early published travel journals describe similar episodes.
In early 20th-century Indonesia: On Salibabu island, North Sulawesi, a 14-year-old boy was killed and supposedly eaten by a specimen in length. Another incident involved a woman reputedly eaten by a "large reticulated python", but few details are known.
In the early 1910s or in 1927, a jeweller went hunting with his friends and was apparently eaten by a python after he sought shelter from a rainstorm in or under a tree. Supposedly, he was swallowed feet-first, perhaps the easiest way for a snake to actually swallow a human.
Among a small group of Aeta peoples in the Philippines, six deaths by pythons were said to have been documented within a period of 40 years, plus one who died later of an infected bite.
In September 1995, a 29-year-old rubber tapper from the southern Malaysian state of Johor was reported to have been killed by a large reticulated python. The victim had apparently been caught unaware and was squeezed to death. The snake had coiled around the lifeless body with the victim's head gripped in its jaws when it was stumbled upon by the victim's brother. The python, reported as measuring long and weighing more than , was killed soon after by the arriving police, who shot it four times.
In October 2003, a woman was reported to be eaten by a giant reticulated python at Sajek Valley in Rangamati Hill District, Bangladesh, when she was collecting paddy crops with her husband. People came to help and retrieved the woman's body from the python's belly.
In October 2008, a woman from Virginia Beach appeared to have been killed by a pet reticulated python. The apparent cause of death was asphyxiation. The snake was later found in the bedroom in an agitated state.
In January 2009, a 3-year-old boy was wrapped in the coils of a pet reticulated python, turning blue. The boy's mother, who had been petsitting the python on behalf of a friend, rescued him by gashing the python with a knife. The snake was later euthanized because of its wounds.
In December 2013, a 59-year-old security guard was strangled to death while trying to capture a python near the Bali Hyatt, a luxury hotel on Indonesia's resort island. The incident happened around 3 am as the 4.5-m (15-ft) python was crossing a road near the hotel. The victim had offered to help capture the snake, which had been spotted several times before near the hotel in the Sanur, Bali, area and escaped back into nearby bushes.
In March 2017, the body of Akbar Salubiro, a 25-year-old farmer in Central Mamuju Regency, West Sulawesi, Indonesia, was found inside the stomach of a reticulated python. He had been declared missing from his palm tree plantation, and the people searching for him found the python the next day with a large bulge in its stomach. They killed the python and found the whole body of the missing farmer inside. This was the first fully confirmed case of a person being eaten by a python. The process of retrieving the body from the python's stomach was documented by pictures and videos taken by witnesses.
In June 2018, a 54-year-old Indonesian woman in Muna Island, Southeast Sulawesi, Indonesia, was killed and eaten by a python. The woman went missing one night while working in her garden, and the next day, a search party was organized after some of her belongings were found abandoned in the garden. The python was found near the garden with a large bulge in its body. The snake was killed and carried into town, where it was cut open, revealing the woman's body completely intact.
In June 2020, a 16-year-old Indonesian boy was attacked and killed by a long python in Bombana Regency, Southeast Sulawesi, Indonesia. The incident took place near a waterfall at Mount Kahar in Rumbia sub-district. The victim was separated from his four friends in the woods. When he screamed, his friends came to help and found him encoiled by a large python. Villagers came to help and managed to kill the snake using a parang machete. However, the victim had already suffocated.
In October 2022, a 52-year-old woman in Terjun Gajah village, Betara Subdistrict, West Tanjung Jabung Regency, Jambi, Indonesia, was killed and swallowed whole by a reticulated python. She went to tap rubber sap on 23 October 2022 and did not return home after sunset. After she was reported missing for a day and a night, a search party discovered a large python with a bulge in its body in a jungle near the rubber plantation. The villagers immediately killed and dissected the python and discovered the intact body of the missing woman inside. Villagers feared more large pythons might be present because farmers previously had reported two goats missing.
In June 2024, a woman of Kalempang village in South Sulawesi province in Indonesia went missing, and her body was discovered inside a reticulated python. 3 weeks later, in July 2024, another woman was discovered inside a python's stomach in South Sulawesi. In August 2024, an elderly woman was found dead after a predatory attempt by a long python. The snake had killed the victim and tried to swallow her, but could not get over the shoulders, regurgitating the body instead. 2 weeks later another woman in Jambi province was killed by a python, which managed to swallow half of her body before being found and killed by the villagers. In November 2024, a 30-year-old man was killed and swallowed whole by a long specimen, the first recorded case of an adult male being eaten since 2017.
In captivity
Increased popularity of the reticulated python in the pet trade is due largely to increased efforts in captive breeding and selectively bred mutations such as the "albino" and "tiger" strains. Other notable color mutations recorded in this species include "sunfire", "motley", "Aztec", "ocelot", "rainbow", and "goldenchild". Some mutations, such as the "BEL" (black eyed leucistic) mutations are seldom bred due the discovery that they develop fatal gastrointestinal issues upon reaching adulthood. The "jaguar" color mutation is likewise controversial in the breeding of captive reticulated pythons, as certain specimens will develop neurological and osseous inner ears, similar to the "Spider" mutation that occurs in the ball python. It remains unclear as to why some individual pythons with the jaguar mutation display neurological issues associated with this genetic trait and others do not.
Smaller animals such as the proposed "super dwarf" subspecies found on small islands are likewise popular due to their smaller size, as they grow to a fraction of the lengths and weights of their mainland kin due to genetics, limited space and prey availability. Dwarf and super dwarf reticulated pythons are likewise defined for captive animals as any reticulated pythons with at least 50 percent lineage hailing from seven select islands in the Selayer island chain near Sulawesi. It can make a good captive, but keepers working with adults from mainland populations should have previous experience with large constrictors to ensure safety to both animal and keeper. Although its interactivity and beauty draws much attention, some feel it is unpredictable. The python can bite and possibly constrict if it feels threatened, or mistakes a hand for food. While not venomous, large pythons can inflict serious injuries by biting, sometimes requiring stitches.
| Biology and health sciences | Snakes | Animals |
88750 | https://en.wikipedia.org/wiki/G-spot | G-spot | The G-spot, also called the Gräfenberg spot (for German gynecologist Ernst Gräfenberg), is characterized as an erogenous area of the vagina that, when stimulated, may lead to strong sexual arousal, powerful orgasms and potential female ejaculation. It is typically reported to be located up the front (anterior) vaginal wall between the vaginal opening and the urethra and is a sensitive area that may be part of the female prostate.
The existence of the G-spot has not been proven, nor has the source of female ejaculation. Although the G-spot has been studied since the 1940s, disagreement persists over its existence as a distinct structure, definition and location. The G-spot may be an extension of the clitoris, which together may be the cause of orgasms experienced vaginally. Sexologists and other researchers are concerned that women may consider themselves to be dysfunctional if they do not experience G-spot stimulation, and emphasize that not experiencing it is normal.
Theorized structure
Location
Two primary methods have been used to define and locate the G-spot as a sensitive area in the vagina: self-reported levels of arousal during stimulation, and stimulation of the G-spot leading to female ejaculation. Ultrasound technology has also been used to identify physiological differences between women, and changes to the G-spot region during sexual activity.
The location of the G-spot is typically reported as being about 50 to 80 mm (2 to 3 in) inside the vagina, on the front wall. For some women, stimulating this area creates a more intense orgasm than clitoral stimulation. The G-spot area has been described as needing direct stimulation, such as two fingers pressed deeply into it. Attempting to stimulate the area through sexual penetration, especially in the missionary position, is difficult because of the particular angle of penetration required.
Vagina and clitoris
Women usually need direct clitoral stimulation in order to orgasm, and G-spot stimulation may be best achieved by using both manual stimulation and vaginal penetration. A yoni massage also includes manual stimulation of the G-spot.
Sex toys are available for G-spot stimulation. One common sex toy is the specially-designed G-spot vibrator, which is a phallic vibrator that has a curved tip and attempts to make G-spot stimulation easy. G-spot vibrators are made from the same materials as regular vibrators, ranging from hard plastic, rubber, silicone, jelly, or any combination of them. The level of vaginal penetration when using a G-spot vibrator depends on the woman, because women's physiology is not always the same. The effects of G-spot stimulation when using the penis or a G-spot vibrator may be enhanced by additionally stimulating other erogenous zones on a woman's body, such as the clitoris or vulva as a whole. When using a G-spot vibrator, this may be done by manually stimulating the clitoris, including by using the vibrator as a clitoral vibrator, or, if the vibrator is designed for it, by applying it so that it stimulates the head of the clitoris, the rest of the vulva and the vagina simultaneously.
A 1981 case study reported that stimulation of the anterior vaginal wall made the area grow by fifty percent and that self-reported levels of arousal/orgasm were deeper when the G-spot was stimulated. Another study, in 1983, examined eleven women by palpating the entire vagina in a clockwise fashion, and reported a specific response to stimulation of the anterior vaginal wall in four of the women, concluding that the area is the G-spot. In a 1990 study, an anonymous questionnaire was distributed to 2,350 professional women in the United States and Canada with a subsequent 55% return rate. Of these respondents, 40% reported having a fluid release (ejaculation) at the moment of orgasm, and 82% of the women who reported the sensitive area (Gräfenberg spot) also reported ejaculation with their orgasms. Several variables were associated with this perceived existence of female ejaculation.
Some research suggests that G-spot and clitoral orgasms are of the same origin. Masters and Johnson were the first to determine that the clitoral structures surround and extend along and within the labia. Upon studying women's sexual response cycle to different stimulation, they observed that both clitoral and vaginal orgasms had the same stages of physical response, and found that the majority of their subjects could only achieve clitoral orgasms, while a minority achieved vaginal orgasms. On this basis, Masters and Johnson argued that clitoral stimulation is the source of both kinds of orgasms, reasoning that the clitoris is stimulated during penetration by friction against its hood.
Researchers at the University of L'Aquila, using ultrasonography, presented evidence that women who experience vaginal orgasms are statistically more likely to have thicker tissue in the anterior vaginal wall. The researchers believe these findings make it possible for women to have a rapid test to confirm whether or not they have a G-spot. Professor of genetic epidemiology, Tim Spector, who co-authored research questioning the existence of the G-spot and finalized it in 2009, also hypothesizes thicker tissue in the G-spot area; he states that this tissue may be part of the clitoris and is not a separate erogenous zone.
Supporting Spector's conclusion is a study published in 2005 which investigates the size of the clitoris – it suggests that clitoral tissue extends into the anterior wall of the vagina. The main researcher of the studies, Australian urologist Helen O'Connell, asserts that this interconnected relationship is the physiological explanation for the conjectured G-spot and experience of vaginal orgasms, taking into account the stimulation of the internal parts of the clitoris during vaginal penetration. While using MRI technology, O'Connell noted a direct relationship between the legs or roots of the clitoris and the erectile tissue of the "clitoral bulbs" and corpora, and the distal urethra and vagina. "The vaginal wall is, in fact, the clitoris," said O'Connell. "If you lift the skin off the vagina on the side walls, you get the bulbs of the clitoris – triangular, crescental masses of erectile tissue." O'Connell et al., who performed dissections on the female genitals of cadavers and used photography to map the structure of nerves in the clitoris, were already aware that the clitoris is more than just its glans and asserted in 1998 that there is more erectile tissue associated with the clitoris than is generally described in anatomical textbooks. They concluded that some females have more extensive clitoral tissues and nerves than others, especially having observed this in young cadavers as compared to elderly ones, and therefore whereas the majority of females can only achieve orgasm by direct stimulation of the external parts of the clitoris, the stimulation of the more generalized tissues of the clitoris via intercourse may be sufficient for others.
French researchers Odile Buisson and Pierre Foldès reported similar findings to those of O'Connell's. In 2008, they published the first complete 3D sonography of the stimulated clitoris, and republished it in 2009 with new research, demonstrating the ways in which erectile tissue of the clitoris engorges and surrounds the vagina. On the basis of this research, they argued that women may be able to achieve vaginal orgasm via stimulation of the G-spot because the highly innervated clitoris is pulled closely to the anterior wall of the vagina when the woman is sexually aroused and during vaginal penetration. They assert that since the front wall of the vagina is inextricably linked with the internal parts of the clitoris, stimulating the vagina without activating the clitoris may be next to impossible. In their 2009 published study, the "coronal planes during perineal contraction and finger penetration demonstrated a close relationship between the root of the clitoris and the anterior vaginal wall". Buisson and Foldès suggested "that the special sensitivity of the lower anterior vaginal wall could be explained by pressure and movement of clitoris's root during a vaginal penetration and subsequent perineal contraction".
Female prostate
In 2001, the Federative Committee on Anatomical Terminology accepted female prostate as a second term for the Skene's gland, which is believed to be found in the G-spot area along the walls of the urethra. The male prostate is biologically homologous to the Skene's gland; it has been unofficially called the male G-spot because it can also be used as an erogenous zone.
Regnier de Graaf, in 1672, observed that the secretions (female ejaculation) by the erogenous zone in the vagina lubricate "in agreeable fashion during coitus". Modern scientific hypotheses linking G-spot sensitivity with female ejaculation led to the idea that non-urine female ejaculate may originate from the Skene's gland, with the Skene's gland and male prostate acting similarly in terms of prostate-specific antigen and prostate-specific acid phosphatase studies, which led to a trend of calling the Skene's glands the female prostate. Additionally, the enzyme PDE5 (involved with erectile dysfunction) has additionally been associated with the G-spot area. Because of these factors, it has been argued that the G-spot is a system of glands and ducts located within the anterior (front) wall of the vagina. A similar approach has linked the G-spot with the urethral sponge.
Clinical significance
G-spot amplification (also called G-spot augmentation or the G-Shot) is a procedure intended to temporarily increase pleasure in sexually active women with normal sexual function, focusing on increasing the size and sensitivity of the G-spot. G-spot amplification is performed by attempting to locate the G-spot and noting measurements for future reference. After numbing the area with a local anesthetic, human engineered collagen is then injected directly under the mucosa in the area the G-spot is concluded to be in.
A position paper published by the American College of Obstetricians and Gynecologists in 2007 warns that there is no valid medical reason to perform the procedure, which is not considered routine or accepted by the College; and it has not been proven to be safe or effective. The potential risks include sexual dysfunction, infection, altered sensation, dyspareunia, adhesions and scarring. The College position is that it is untenable to recommend the procedure. The procedure is also not approved by the Food and Drug Administration or the American Medical Association, and no peer-reviewed studies have been accepted to account for either safety or effectiveness of this treatment.
Society and culture
General skepticism
In addition to general skepticism among gynecologists, sexologists and other researchers that the G-spot exists, a team at King's College London in late 2009 suggested that its existence is subjective. They acquired the largest sample size of women to date – 1,800 – who are pairs of twins, and found that the twins did not report a similar G-spot in their questionnaires. The research, headed by Tim Spector, documents a 15-year study of the twins, identical and non-identical. According to the researchers, if one identical twin reported having a G-spot, it was more likely that the other would too, but this pattern did not materialize. Study co-author Andrea Burri believes: "It is rather irresponsible to claim the existence of an entity that has never been proven and pressurise women and men too." She stated that one of the reasons for the research was to remove feelings of "inadequacy or underachievement" for women who feared they lacked a G-spot. Researcher Beverly Whipple dismissed the findings, commenting that twins have different sexual partners and techniques, and that the study did not properly account for lesbian or bisexual women.
Petra Boynton, a British scientist who has written extensively on the G-spot debate, is also concerned about the promotion of the G-spot leading women to feel "dysfunctional" if they do not experience it. "We're all different. Some women will have a certain area within the vagina which will be very sensitive, and some won't — but they won't necessarily be in the area called the G spot," she stated. "If a woman spends all her time worrying about whether she is normal, or has a G spot or not, she will focus on just one area, and ignore everything else. It's telling people that there is a single, best way to have sex, which isn't the right thing to do."
Nerve endings
G-spot proponents are criticized for giving too much credence to anecdotal evidence, and for questionable investigative methods; for instance, the studies which have yielded positive evidence for a precisely located G-spot involve small participant samples. While the existence of a greater concentration of nerve endings at the lower third (near the entrance) of the vagina is commonly cited, some scientific examinations of vaginal wall innervation have shown no single area with a greater density of nerve endings.
Several researchers also consider the connection between the Skene's gland and the G-spot to be weak. The urethral sponge, however, which is also hypothesized as the G-spot, contains sensitive nerve endings and erectile tissue. Sensitivity is not determined by neuron density alone: other factors include the branching patterns of neuron terminals and cross or collateral innervation of neurons. While G-spot opponents argue that because there are very few tactile nerve endings in the vagina and that therefore the G-spot cannot exist, G-spot proponents argue that vaginal orgasms rely on pressure-sensitive nerves.
Clitoral and other anatomical debates
The G-spot having an anatomical relationship with the clitoris has been challenged by Vincenzo Puppo, who, while agreeing that the clitoris is the center of female sexual pleasure, disagrees with Helen O'Connell and other researchers' terminological and anatomical descriptions of the clitoris. He stated, "Clitoral bulbs is an incorrect term from an embryological and anatomical viewpoint, in fact the bulbs do not develop from the phallus, and they do not belong to the clitoris". He says that clitoral bulbs "is not a term used in human anatomy" and that vestibular bulbs is the correct term, adding that gynecologists and sexual experts should inform the public with facts instead of hypotheses or personal opinions. "[C]litoral/vaginal/uterine orgasm, G/A/C/U spot orgasm, and female ejaculation, are terms that should not be used by sexologists, women, and mass media", he said, further commenting that the "anterior vaginal wall is separated from the posterior urethral wall by the urethrovaginal septum (its thickness is 10–12 mm)" and that the "inner clitoris" does not exist. "The female perineal urethra, which is located in front of the anterior vaginal wall, is about one centimeter in length and the G-spot is located in the pelvic wall of the urethra, 2–3 cm into the vagina", Puppo stated. He believes that the penis cannot come in contact with the congregation of multiple nerves/veins situated until the angle of the clitoris, detailed by Georg Ludwig Kobelt, or with the roots of the clitoris, which do not have sensory receptors or erogenous sensitivity, during vaginal intercourse. He did, however, dismiss the orgasmic definition of the G-spot that emerged after Ernst Gräfenberg, stating that "there is no anatomical evidence of the vaginal orgasm which was invented by Freud in 1905, without any scientific basis".
Puppo's belief that there is no anatomical relationship between the vagina and clitoris is contrasted by the general belief among researchers that vaginal orgasms are the result of clitoral stimulation; they maintain that clitoral tissue extends, or is at least likely stimulated by the clitoral bulbs, even in the area most commonly reported to be the G-spot. "My view is that the G-spot is really just the extension of the clitoris on the inside of the vagina, analogous to the base of the male penis", said researcher Amichai Kilchevsky. Because female fetal development is the "default" direction of fetal development in the absence of substantial exposure to male hormones and therefore the penis is essentially a clitoris enlarged by such hormones, Kilchevsky believes that there is no evolutionary reason why females would have two separate structures capable of producing orgasms and blames the porn industry and "G-spot promoters" for "encouraging the myth" of a distinct G-spot.
The general difficulty of achieving vaginal orgasms, which is a predicament that is likely due to nature easing the process of childbearing by drastically reducing the number of vaginal nerve endings, challenge arguments that vaginal orgasms help encourage sexual intercourse in order to facilitate reproduction. O'Connell stated that focusing on the G-spot to the exclusion of the rest of a woman's body is "a bit like stimulating a guy's testicles without touching the penis and expecting an orgasm to occur just because love is present". She stated that it "is best to think of the clitoris, urethra, and vagina as one unit because they are intimately related". Ian Kerner stated that the G-spot may be "nothing more than the roots of the clitoris crisscrossing the urethral sponge".
A Rutgers University study, published in 2011, was the first to map the female genitals onto the sensory portion of the brain, and supports the possibility of a distinct G-spot. When the research team asked several women to stimulate themselves in a functional magnetic resonance (fMRI) machine, brain scans showed stimulating the clitoris, vagina and cervix lit up distinct areas of the women's sensory cortex, which means the brain registered distinct feelings between stimulating the clitoris, the cervix and the vaginal wall – where the G-spot is reported to be. "I think that the bulk of the evidence shows that the G-spot is not a particular thing," stated Barry Komisaruk, head of the research findings. "It's not like saying, 'What is the thyroid gland?' The G-spot is more of a thing like New York City is a thing. It's a region, it's a convergence of many different structures".
In 2009, The Journal of Sexual Medicine held a debate for both sides of the G-spot issue, concluding that further evidence is needed to validate the existence of the G-spot. In 2012, scholars Kilchevsky, Vardi, Lowenstein and Gruenwald stated in the journal, "Reports in the public media would lead one to believe the G-spot is a well-characterized entity capable of providing extreme sexual stimulation, yet this is far from the truth". The authors cited that dozens of trials have attempted to confirm the existence of a G-spot using surveys, pathologic specimens, various imaging modalities, and biochemical markers, and concluded:The surveys found that a majority of women believe a G-spot actually exists, although not all of the women who believed in it were able to locate it. Attempts to characterize vaginal innervation have shown some differences in nerve distribution across the vagina, although the findings have not proven to be universally reproducible. Furthermore, radiographic studies have been unable to demonstrate a unique entity, other than the clitoris, whose direct stimulation leads to vaginal orgasm. Objective measures have failed to provide strong and consistent evidence for the existence of an anatomical site that could be related to the famed G-spot. However, reliable reports and anecdotal testimonials of the existence of a highly sensitive area in the distal anterior vaginal wall raise the question of whether enough investigative modalities have been implemented in the search of the G-spot.
A 2014 review from Nature Reviews Urology reported that "no single structure consistent with a distinct G-spot has been identified".
History
The release of fluids had been seen by medical practitioners as beneficial to health. Within this context, various methods were used over the centuries to release "female seed" (via vaginal lubrication or female ejaculation) as a treatment for suffocation ex semine retento (suffocation of the womb), female hysteria or green sickness. Methods included a midwife rubbing the walls of the vagina or insertion of the penis or penis-shaped objects into the vagina. In the book History of V, lists old terms for what she believes refer to the female prostate (the Skene's gland), including the little stream, the black pearl and palace of yin in China, the skin of the earthworm in Japan, and saspanda nadi in the India sex manual Ananga Ranga.
The 17th-century Dutch physician Regnier de Graaf described female ejaculation and referred to an erogenous zone in the vagina that he linked as homologous with the male prostate; this zone was later reported by the German gynecologist Ernst Gräfenberg. Coinage of the term G-spot has been credited to Addiego et al. in 1981, named after Gräfenberg, and to Alice Kahn Ladas and Beverly Whipple et al. in 1982. Gräfenberg's 1940s research, however, was dedicated to urethral stimulation; Gräfenberg stated, "An erotic zone always could be demonstrated on the anterior wall of the vagina along the course of the urethra". The concept of the G-spot entered popular culture with the 1982 publication of The G Spot and Other Recent Discoveries About Human Sexuality by Ladas, Whipple and Perry, but it was criticized immediately by gynecologists: some of them denied its existence as the absence of arousal made it less likely to observe, and autopsy studies did not report it.
| Biology and health sciences | Human anatomy | Health |
88793 | https://en.wikipedia.org/wiki/Eagle%20ray | Eagle ray | The eagle rays are a group of cartilaginous fishes in the family Myliobatidae, consisting mostly of large species living in the open ocean rather than on the sea bottom.
Eagle rays feed on mollusks, and crustaceans, crushing their shells with their flattened teeth. They are excellent swimmers and are able to breach the water up to several meters above the surface. Compared with other rays, they have long tails, and well-defined, rhomboidal bodies. They are ovoviviparous, giving birth to up to six young at a time. They range from in length and 7 m (23 ft) in wingspan.
Classification
Nelson's book Fishes of the World treats cownose rays, mantas, and devil rays as subfamilies in the Myliobatidae. However, most authors (including William Toby White) have preferred to leave the Rhinopteridae and Mobulidae outside of the Myliobatidae. White (2014) retained three genera (Aetobatus, Aetomylaeus, and Myliobatis) in the Myliobatidae, while a fourth (Pteromylaeus) was synonymized with Aetomylaeus. A 2016 paper placed Aetobatus in its own family, the Aetobatidae.
| Biology and health sciences | Batoidea | Animals |
88837 | https://en.wikipedia.org/wiki/Nile%20perch | Nile perch | The Nile perch (Lates niloticus), also known as the African snook, Goliath perch, African barramundi, Goliath barramundi, Giant lates or the Victoria perch, is a species of freshwater fish in family Latidae of order Perciformes. It is widespread throughout much of the Afrotropical realm, being native to the Congo, Nile, Senegal, Niger and Lake Chad, Volta, Lake Turkana, and other river basins. It also occurs in the brackish waters of Lake Maryut in Egypt. The Nile perch is a fish of substantial economic and food-security importance in East Africa. Originally described as Labrus niloticus, among the marine wrasses, the species has also been referred to as Centropomus niloticus. Common names include African snook, Victoria perch (a misleading trade name, as the species is not native to Lake Victoria, though they have been introduced there), and many local names in various African languages, such as the Luo name mbuta or mputa. In Tanzania, it is called sangara, sankara, or chenku. In Francophone African countries, it is known as capitaine. Its name in the Hausa language is giwan ruwa, meaning "water elephant".
Description
L. niloticus is silver in color with a blue tinge. It has distinctive dark-black eyes, with a bright-yellow outer ring. One of the largest freshwater fishes, it reaches a maximum length of nearly , weighing up to . Mature fish typically range from , although many fish are caught before they can grow this large.
Adult Nile perch occupy all habitats of a lake with sufficient oxygen concentrations, while juveniles are restricted to shallow or nearshore environments. A fierce predator that dominates its surroundings, the Nile perch feeds on fish (including its own species), crustaceans, molluscs, and insects; the juveniles also feed on zooplankton. Nile perch use schooling as a mechanism to protect themselves from other predators.
Invasive species
Nile perch have been introduced to many other lakes in Africa, including Lake Victoria and the artificial Lake Nasser. The World Conservation Union's Invasive Species Specialist Group considers L. niloticus one of the world's 100 worst invasive species.
The state of Queensland in Australia levies heavy fines on anyone found in possession of a living Nile perch, since it competes directly with the native barramundi, which is similar and grows to long, while the Nile perch grows to long.
The species is of great commercial importance as a food fish. The Nile perch is also popular with sport anglers, as it attacks artificial fishing lures, and it is also raised in aquaculture.
Lake Victoria introduction
The introduction of this species to Lake Victoria is one of the most cited examples of the negative effects alien species can have on ecosystems.
The Nile perch was introduced to Lake Victoria in East Africa in the 1950s, and has since been fished commercially. In 2003, Nile perch sales to the EU reached 169 million euros. Sport-fishing in the region of Uganda and Tanzania provided additional income from tourism.
Its introduction was ecologically disruptive and is attributed with causing the extinction or near-extinction of several hundred native species, with some populations fluctuating with commercial fishing and the actual Nile perch stocks. The Nile perch initially fed on native cichlids, but with decreasing availability of this prey, it now consumes mainly small shrimp and minnows.
The alteration of the native ecosystem had disruptive socioeconomic effects on local communities bordering the lake. Many local people have been displaced from their traditional occupations in the fishing trade and brought them into the cash economy, or before the establishment of export-oriented fisheries, turned them into economic refugees. At least initially, nets strong enough to hold adult Nile perch could not be manufactured locally and had to be imported for a high price.
The introduction of Nile perch has also had additional ecological effects on shore. Native cichlids were traditionally sun-dried, but because Nile perch have a high fat content (higher than cichlids), they need to be smoked to avoid spoiling. This has led to an increased demand for firewood in a region already hard-hit by deforestation, soil erosion, and desertification.
The Academy Award-nominated documentary Darwin's Nightmare by Hubert Sauper (a French-Austrian-Belgian production, 2004) deals with the damage that has been caused by Nile perch introduction, including the import of weapons and ammunition in cargo planes from Europe, which are then used to export Nile perch, further exacerbating conflict and misery in the surrounding regions.
Regardless of whether it is considered positive or negative, the trophic web of Lake Victoria appears to have been drastically impoverished by the introduction of this novel near-top-level predator. While the ecosystem seems to be moving towards a new equilibrium, neither its former state nor the state of fisheries on Lake Victoria can ever easily be brought back.
Threats
Despite being a successful invasive species, the fish faces threats. Being a species of megafauna, the most obvious threats to the species are overfishing and the use of illegal fishing gear, as well as invasive water hyacinths. Prey depletion is also a factor, as it decreases the size of the fish and makes it vulnerable to larger predators, such as crocodiles.
Export
In 2021, The Uganda Fish Processors and Exporters Association called on the parliament to ban the local consumption of the species so as to protect its export.
Religion
Nile perch were involved with the worship of Neith. As a result nile perch were sometimes mummified. A deposit of several thousand mummified perch was excavated in an area to the west of Esna where there was a temple to Neith. Mummified perch have also been found at Gurob near a temple to Neith while perch statuettes have been found at Sais again in the context of a temple to Neith.
| Biology and health sciences | Acanthomorpha | Animals |
88847 | https://en.wikipedia.org/wiki/Philippine%20eagle | Philippine eagle | The Philippine eagle (Pithecophaga jefferyi), also known as the monkey-eating eagle or great Philippine eagle, is a critically endangered species of eagle of the family Accipitridae which is endemic to forests in the Philippines. It has brown and white-colored plumage, a shaggy crest, and generally measures in length and weighs .
The Philippine eagle is considered the largest of the extant eagles in the world in terms of length and wing surface area, with only Steller's sea eagle and the Harpy eagle being larger in terms of weight and bulk. It has been declared the national bird of the Philippines. The most significant threat to the species is loss of habitat, a result of high levels of deforestation throughout most of its range. Because of this, the IUCN Red List has classified the species as "critically endangered".
Killing a Philippine eagle is a criminal offence, punishable by law with up to 12 years' imprisonment and heavy fines.
Names
The Philippine eagle has numerous native names in the Philippine languages. These include and (a Spanish loanword) in Tagalog; or in the Visayan languages; or in Maranao and Maguindanao; in Subanen; and or in the Manobo languages, Klata, Tagabawa, Mandaya, and Kalagan. Some of these names are also used for other large eagles in general, such as the white-breasted sea eagle (also called in Visayan). In modern Filipino, it is usually referred to as háribon (a portmanteau of "haring ibón", literally "bird king").
Taxonomy
The first European to study the species was the English explorer and naturalist John Whitehead in 1896, who observed the bird and whose servant, Juan, collected the first specimen a few weeks later. The skin of the bird was sent to William Robert Ogilvie-Grant in London in 1896, who initially showed it off in a local restaurant and described the species a few weeks later.
Upon its scientific discovery, the Philippine eagle was first called the monkey-eating eagle because of reports from natives of Bonga, Samar, where the species was first discovered, that it preyed exclusively on monkeys. These reports gave its generic name, from the Greek (πίθηκος, "ape" or "monkey") and (-φάγος, "eater of"). The species name commemorates Jeffery Whitehead, the father of John Whitehead. Later studies revealed, however, that the alleged monkey-eating eagle also ate other animals, such as colugos, large snakes, monitor lizards, and even large birds such as Hornbills. This, coupled with the fact that the same name applied to the African Crowned eagle and the Central and South American Harpy eagle, it was renamed "Philippine eagle" in a 1978 proclamation by then-President Ferdinand Marcos. In 1995, it was declared a national emblem under President Fidel V. Ramos. This species has no recognized subspecies.
Evolutionary history
A 1919 study of the bird’s skeletal features led to the suggestion that the nearest relative was the Harpy eagle (Harpia harpyja). The species was included in the subfamily Harpiinae until a 2005 study of DNA sequences which identified them as not members of the group, finding instead that the nearest relatives are Snake eagles (Circaetinae), such as the bateleur. The species has subsequently been placed in the subfamily Circaetinae.
Description
The Philippine eagle's nape is adorned with long, brown feathers that form a shaggy, mane-like crest. The eagle has a dark face and a creamy-brown nape and crown. The back of the Philippine eagle is dark brown, while the underside and underwings are white. The heavy legs are yellow, with large, powerful, dark claws, and the prominent, large, high-arched, deep beak is a bluish-gray. The eagle's eyes are blue-gray. Juveniles are similar to adults except their upperpart feathers have pale fringes.
The Philippine eagle is typically reported as measuring in total length, but a survey of several specimens from some of the largest natural history collections in the world found the average was for males and for females. Based on the latter measurements, this makes it the longest extant species of eagle, as the average for the female equals the maximum reported for the harpy eagle and Steller's sea eagle. The longest Philippine eagle reported anywhere and the longest eagle outside of the extinct Haast's eagle is a specimen from Field Museum of Natural History with a length of , but it had been kept in captivity so may not represent the wild individuals due to differences in the food availability.
The level of sexual dimorphism in size is not certain, but the male is believed to be typically about 10% smaller than the female, and this is supported by the average length provided for males and females in one source. In many of the other large eagle species, the size difference between adult females and males can exceed 20%. For adult Philippine eagles, the complete weight range has been reported as , while others have found the average was somewhat lower than the above range would indicate, at for males and for females. One male (age not specified) was found to weigh . The Philippine eagle has a wingspan of and a wing chord length of .
The maximum reported weight is surpassed by two other eagles (the harpy and Steller's sea eagles) and the wings are shorter than large eagles of open country (such as the white-tailed eagle, Steller's sea eagle, martial eagle, or wedge-tailed eagle), but are quite broad. The tarsus of the Philippine eagle ties as the longest of any eagle from long, which is about the same length as that of the much smaller but relatively long-legged New Guinea eagle. The very large but laterally compressed bill rivals the size of Steller's sea eagle's as the largest bill for an extant eagle. Its bill averages in length from the gape. The tail is fairly long at , while another source lists a tail length of .
The most frequently heard noises made by the Philippine eagle are loud, high-pitched whistles ending with inflections in pitch. Additionally, juveniles have been known to beg for food by a series of high-pitched calls.
Distribution and habitat
The Philippine eagle is endemic to the Philippines and can be found on four major islands: eastern Luzon, Samar, Leyte, and Mindanao. The largest numbers of eagles reside on Mindanao, with between 82 and 233 breeding pairs. Only six pairs are found on Samar, two on Leyte, and a few on Luzon. It can be found in Northern Sierra Madre National Park on Luzon and Mount Apo, Mount Malindang, and Mount Kitanglad National Parks on Mindanao.
This eagle is found in dipterocarp and mid montane forests, particularly in steep areas. Its elevation ranges from the lowlands to mountains of over . Only an estimated of old-growth forest remain in the bird's range. However, its total estimated range is about .
Ecology and behavior
Evolution in the Philippines, without other predators, made the eagles the dominant hunter in the Philippine forests. The Philippine eagle has a wide range of prey which includes birds, reptiles and mammals (mainly civets and colugos).
Each breeding pair requires a large home range to successfully raise a chick, thus the species is extremely vulnerable to deforestation. Earlier, the territory has been estimated at , but a study on Mindanao Island found the nearest distance between breeding pairs to be about on average, resulting in a circular plot of .
The species' flight is fast and agile, resembling the smaller hawks more than similar large birds of prey.
Juveniles in play behavior have been observed gripping knotholes in trees with their talons, and using their tails and wings for balance, inserting their heads into tree cavities. Additionally, they have been known to attack inanimate objects for practice, as well as attempt to hang upside down to work on their balance. As the parents are not nearby when this occurs, they apparently do not play a role in teaching the juvenile to hunt.
Life expectancy for a wild eagle is estimated to be from 30 to 60 years. A captive Philippine eagle lived for 41 years in Rome Zoo, and it was already an adult when it arrived at the zoo. Another captive Philippine eagle lived for 46 years at the Philippine Eagle Center in Davao City. However, wild birds on average are believed to live shorter lives than captive birds.
Diet
The Philippine eagle was known initially as the Philippine monkey-eating eagle because it was believed to feed on monkeys almost exclusively. The only two monkeys native to the Philippines are Philippine long-tailed macaque and common long-tailed macaques, both are subspecies of crab-eating macaque and weighing in males and in females. Though Philippine eagles do prey on these monkeys, they are an opportunist apex predator, taking prey based on their local level of abundance and ease. This misconception may have come from the first examined specimen which was found to have undigested pieces of a monkey in its stomach.
Prey specimens found at the eagle's nest have ranged in size from a small bat weighing to a Philippine deer weighing . The primary prey is usually the tree squirrel-sized Philippine flying lemurs, which can make up an estimated 90% of the raptor's diet in some locations. However, primary prey species vary from island to island depending on species availability, particularly in Luzon and Mindanao, because the islands are in different faunal regions. For example, flying lemurs are preferred prey in Mindanao, but are absent in Luzon. The primary prey for the eagles seen in Luzon are macaques, reptiles, and Northern Luzon giant cloud rat (Phloeomys pallidus) which can weigh twice as much as flying lemurs at . In many regions, civets are taken as supplemental prey, mainly cat-sized Asian palm civets (Paradoxurus hermaphroditus) but occasionally larger Malay civets (Viverra tangalunga). Other mammalian prey can include flying squirrels, tree squirrels, flying foxes, rats and mouse-deer. Birds are also taken, including large species such as owls, hawks and hornbills. Reptiles occasionally form a large part of their diet, snakes such as Ptyas luzonensis and Gonyosoma oxycephalum are mainly taken. Venomous pit vipers (Viperidae) are also taken as prey, and in one instance, a breeding pair delivered a Philippine cobra (Naja philippinensis) to the nest. Monitor lizards, including marbled water monitors (Varanus marmoratus) and even larger Northern Sierra Madre forest monitors (Varanus bitatawa) have been taken. While most of the prey consists of wild prey, they have been reported to capture domestic fowls (Gallus gallus domesticus), cats (Felis catus), young pigs (Sus domesticus) and small dogs (Canis familiaris).
Philippine eagles primarily use two hunting techniques. One is still-hunting, in which it watches for prey activity while sitting almost motionlessly on a branch near the canopy. The other is perch-hunting, which entails periodically gliding from one perch to another. While perch-hunting, they often work their way gradually down from the canopy down the branches, and if not successful in finding prey in their initial foray, they fly or circle back up to the top of the trees to work them again. Eagles in Mindanao often find success using the latter method while hunting flying lemurs, since they are nocturnal animals that try to use camouflage to protect themselves by day. Eagle pairs sometimes hunt troops of monkeys cooperatively, with one bird perching nearby to distract the primates, allowing the other to swoop in from behind, hopefully unnoticed, for the kill. Since the native macaque is aggressive and often around the same size as the eagle itself or even larger, up to in adult males, it is a potentially hazardous prey, and an eagle has been reported to suffer a broken leg after it struggled and fell along with a large male monkey.
Reproduction
The complete breeding cycle of the Philippine eagle lasts two years. The female matures sexually at five years of age and the male at seven. Like most eagles, the Philippine eagle is monogamous. Once paired, a couple remains together for the rest of their lives. If one dies, the remaining eagle often searches for a new mate to replace the one lost.
The beginning of courtship is signaled by nest-building, and the eagle remaining near its nest. Aerial displays also play a major role in the courtship. These displays include paired soaring over a nesting territory, the male chasing the female in a diagonal dive, and mutual talon presentation, where the male presents his talons to the female's back and she flips over in midair to present her own talons. Advertisement displays coupled with loud calling have also been reported. The willingness of an eagle to breed is displayed by the eagle bringing nesting materials to the bird's nest. Copulation follows and occurs repeatedly both on the nest and on nearby perches. The earliest courtship has been reported in July.
Breeding season is in July; birds on different islands, most notably Mindanao and Luzon, begin breeding at different ends of this range. The amount of rainfall and population of prey may also affect the breeding season. The nest is normally built on an emergent dipterocarp, or any tall tree with an open crown, in primary or disturbed forest. The nests are lined with green leaves, and can be around across. The nesting location is around or even more above the ground. As in many other large raptors, the eagle's nest resembles a huge platform made of sticks. The eagle frequently reuses the same nesting site for several different chicks. Eight to 10 days before the egg is ready to be laid, the female is afflicted with a condition known as egg lethargy. In this experience, the female does not eat, drinks much water, and holds her wings droopingly. The female typically lays one egg in the late afternoon or at dusk, although occasionally two have been reported. If an egg fails to hatch or the chick dies early, the parents likely lay another egg the following year. Copulation may take place a few days after the egg is laid to enable another egg to be laid should the first one fail. The egg is incubated for 58 to 68 days (typically 62 days) after being laid. Both sexes participate in the incubation, but the female does the majority of incubating during the day and all of it at night.
Both sexes help feed the newly hatched eaglet. Additionally, the parents have been observed taking turns shielding the eaglet from the sun and rain until it is seven weeks old. The young eaglet fledges after four or five months. The earliest an eagle has been observed making a kill is 304 days after hatching. Both parents take care of the eaglet for a total of 20 months and, unless the previous nesting attempt had failed, the eagles can breed only in alternate years. The Philippine eagle rivals two other large tropical eagles, namely the crowned eagle and harpy eagle, for having the longest breeding cycle of any bird of prey. Even nests have no predators other than humans, as even known nest predators such as palm civets and macaques (being prey species) are likely to actively avoid any area with regular eagle activity.
Conservation
In 1994, the IUCN and BirdLife International listed this species as critically endangered. The IUCN believed that between 180 and 500 Philippine eagles survive in the Philippines. In 2015, about 600 were estimated to be left in the wild. The Zoological Society of London listed the Philippine eagle as the top 15 EDGE species, out of all the world's recorded species, making it the Philippines' most "evolutionary distinct and globally endangered" species.
They are threatened primarily by deforestation through logging and expanding agriculture. Old-growth forest is being lost at a high rate, and most of the forest in the lowlands is owned by logging companies. Mining, pollution, exposure to pesticides that affect breeding, and poaching are also major threats. Additionally, they are occasionally caught in traps laid by local people for deer. Though this is no longer a major problem, the eagle's numbers were also reduced by being captured for zoos. The diminishing numbers of the Philippine eagle were first brought to international attention in 1965 by the noted Filipino ornithologist Dioscoro S. Rabor, and the director of the Parks and Wildlife Office, Jesus A. Alvarez. Charles Lindbergh, best known for crossing the Atlantic alone and without stopping in 1927, was fascinated by this eagle. As a representative of the World Wildlife Fund, Lindbergh traveled to the Philippines several times between 1969 and 1972, where he helped persuade the government to protect the eagle. In 1969, the Monkey-eating Eagle Conservation Program was started to help preserve this species. In 1992, the first Philippine eagles were hatched in captivity through artificial insemination; however, the first naturally bred eaglet was not hatched until 1999. The first captive-bred bird to be released in the wild, Kabayan, was released in 2004 on Mindanao; however, he was accidentally electrocuted in January 2005. Another eagle, Kagsabua, was released in March 2008, but was shot and eaten by a farmer. In June 2015, an eagle was released after being treated for a gunshot wound; two months later it was shot and killed. Killing this critically endangered species is punishable under Philippine law by 12 years in jail and heavy fines.
Its numbers have slowly dwindled over the decades to the current population of 180 to 600 eagles. A series of floods and mud slides, caused by deforestation, further devastated the remaining population. The Philippine eagle may soon no longer be found in the wild, unless direct intervention is taken. The Philippine Eagle Foundation in Davao City, is one organization dedicated to the protection and conservation of the Philippine eagle and its forest habitat. The Philippine Eagle Foundation has successfully bred Philippine eagles in captivity for over a decade and conducted the first experimental release of a captive-bred eagle to the wild. The foundation has 36 eagles at its center, of which 19 were bred in captivity, including one born in 2015, the first for two years.
Ongoing research on behavior, ecology, and population dynamics is also underway. In recent years, protected lands have been established specifically for this species, such as the of Cabuaya Forest and the of Taft Forest Wildlife Sanctuary on Samar. However, a large proportion of the population is found on unprotected land.
Philippine Eagle Diplomacy
In a June 2019 wildlife loan agreement, a pair of Philippine eagles; (as of 2019) Geothermica, a 15 year old male and Sambisig (meaning one unity) a 17 year old female, were loaned to the Jurong Bird Park Singapore. This was part of conservation efforts for the species and also commemorated 50 years friendship between the Philippines and Singapore. Before the pair were shipped, they were given 2 specially printed passports dedicated to them. The pair were scheduled to stay at Singapore for 10 years as part of the agreement. The move marked the initiation of the government's Philippine Eagle Diplomacy program. As of 2024, Geothermica had died from lung infection, his body now taxidermied at the National Museum, while Sambisig remained in Singapore at the new bird park that replaced Jurong.
In August 2024, the Philippine Eagle Foundation launched the Philippine Eagle Geothermica’s Gateway Project facility as tribute to 19-year-old Geothermica, who died from an Aspergillus molds' fungal infection at the Bird Paradise, Mandai Wild Reserve in September 2022. His sculpture stands in the center of the new facility while his taxidermied body is currently displayed at the National Museum of the Philippines. The PEF partnered with Mandai Wildlife Group, United Architects of the Philippines-Davao City and Kublai Millan for the project.
Representation
The Philippine eagle was officially declared the national bird of the Philippines on July 4, 1995, by President Fidel V. Ramos under Proclamation No. 615, series of 1995. Due to the eagle’s size and rarity, it is also a coveted species for birdwatchers.
The Philippine eagle has been featured on at least 12 stamps from the Philippines, with dates ranging from 1967 to 2007. It was also depicted on the 50-sentimo coins minted from 1983 to 1994; in 2018, on the 500-Piso commemorative silver coin, to celebrate the 25th anniversary of Bangko Sentral ng Pilipinas; on January 18, 2021, and on the 5,000-Piso Lapulapu commemorative non-circulating banknote. On December 11, 2021, the Bangko Sentral ng Pilipinas unveiled the design of the new 1,000-Piso polymer banknote featuring the bird as the main portrait, controversially replacing those of three World War II martyrs: Chief Justice José Abad Santos, suffragist Josefa Llanes Escoda, and General Vicente Lim.
Historically, about 50 Philippine eagles have been kept in zoos in Europe (England, Germany, Belgium, Italy and France), the United States, and Japan. The first was a female that arrived in London Zoo in August 1909 and died there in February 1910. The majority arrived in zoos between 1947 and 1965. The last outside the Philippines died in 1988 in the Antwerp Zoo, where it had lived since 1964 (except for a period at the Planckendael Zoo in Belgium). The first captive breeding was only achieved in 1992 at the facility of the Philippine Eagle Foundation in Davao City, Philippines, which has bred it several times since then.
The Philippine eagle is also used in sporting events as a mascot, most notably one in the 2005 Southeast Asian Games held in Manila known as "Gilas". The Philippine eagle is also the animal used in the Philippines men's national basketball team or Gilas Pilipinas' logo/team crest.
| Biology and health sciences | Accipitrimorphae | Animals |
88857 | https://en.wikipedia.org/wiki/Blood%20transfusion | Blood transfusion | Blood transfusion is the process of transferring blood products into a person's circulation intravenously. Transfusions are used for various medical conditions to replace lost components of the blood. Early transfusions used whole blood, but modern medical practice commonly uses only components of the blood, such as red blood cells, plasma, platelets, and other clotting factors. White blood cells are transfused only in very rare circumstances, since granulocyte transfusion has limited applications. Whole blood has come back into use in the trauma setting.
Red blood cells (RBC) contain hemoglobin and supply the cells of the body with oxygen. White blood cells are not commonly used during transfusions, but they are part of the immune system and also fight infections. Plasma is the "yellowish" liquid part of blood, which acts as a buffer and contains proteins and other important substances needed for the body's overall health. Platelets are involved in blood clotting, preventing the body from bleeding. Before these components were known, doctors believed that blood was homogeneous. Because of this scientific misunderstanding, many patients died because of incompatible blood transferred to them.
Medical uses
Red cell transfusion
Historically, red blood cell transfusion was considered when the hemoglobin level fell below 100g/L or hematocrit fell below 30%. Because each unit of blood given carries risks, a trigger level lower than that, at 70 to 80g/L, is now usually used, as it has been shown to have better patient outcomes. The administration of a single unit of blood is the standard for hospitalized people who are not bleeding, with this treatment followed with re-assessment and consideration of symptoms and hemoglobin concentration. Patients with poor oxygen saturation may need more blood. The advisory caution to use blood transfusion only with more severe anemia is in part due to evidence that outcomes are worsened if larger amounts are given. One may consider transfusion for people with symptoms of cardiovascular disease such as chest pain or shortness of breath. In cases where patients have low levels of hemoglobin due to iron deficiency, but are cardiovascularly stable, oral or parenteral iron is a preferred option based on both efficacy and safety. Other blood products are given where appropriate, e.g., fresh frozen plasma to treat clotting deficiencies and platelets to treat or prevent bleeding in thrombocytopenic patients.
Procedure
Before a blood transfusion is given, there are many steps taken to ensure quality of the blood products, compatibility, and safety to the recipient. In 2012, a national blood policy was in place in 70% of countries and 69% of countries had specific legislation that covers the safety and quality of blood transfusion.
Blood donation
The source of blood to be transfused can either be the potential recipient (autologous transfusion), or someone else (allogeneic or homologous transfusion). The latter is much more common than the former. Using another's blood must first start with donation of blood. Blood is most commonly donated as whole blood obtained intravenously and mixed with an anticoagulant. In first-world countries, donations are usually anonymous to the recipient, but products in a blood bank are always individually traceable through the whole cycle of donation, testing, separation into components, storage, and administration to the recipient. This enables management and investigation of any suspected transfusion related disease transmission or transfusion reaction. Developing countries rely heavily on replacement and remunerated donors rather than voluntary nonremunerated donors due to concerns regarding donation- and transfusion-transmitted infection as well as local and cultural beliefs.
It is unclear whether applying alcohol swab alone or alcohol swab followed by antiseptic is able to reduce contamination of donor's blood.
Studies show that the main motivators to blood donation tend to be prosocial (e.g., altruism, selflessness, charity), while the main deterrents include fear, distrust, or perceived racial discrimination in historic contexts.
Processing and testing
Donated blood is usually subjected to processing after it is collected, to make it suitable for use in specific patient populations. Collected blood is then separated into blood components by centrifugation: red blood cells, plasma, and platelets. Plasma can be further processed to manufacture albumin protein, clotting factor concentrates, cryoprecipitate, fibrinogen concentrate, and immunoglobulins (antibodies). Red cells, plasma and platelets can also be donated individually via a more complex process called apheresis.
The World Health Organization (WHO) recommends that all donated blood be tested for transfusion-transmissible infections. These include HIV, hepatitis B, hepatitis C, Treponema pallidum (syphilis) and, where relevant, other infections that pose a risk to the safety of the blood supply, such as Trypanosoma cruzi (Chagas disease) and Plasmodium species (malaria). According to the WHO, 10 countries are not able to screen all donated blood for one or more of: HIV, hepatitis B, hepatitis C, or syphilis. One of the main reasons for this is because testing kits are not always available. However the prevalence of transfusion-transmitted infections is much higher in low income countries compared to middle and high income countries.
All donated blood should also be tested for the ABO blood group system and Rh blood group system to ensure that the patient is receiving compatible blood.
In addition, in some countries platelet products are also tested for bacterial infections due to its higher inclination for contamination due to storage at room temperature.
Donors may be tested for cytomegalovirus (CMV) because of the risk of transmission to certain immunocompromised recipients, such as those with stem cell transplant or T cell diseases. However, testing is not universally mandated, because leukoreduced blood is generally considered safe from CMV transmission; also, most donors (and recipients) are seropositive for CMV, and are not actively viremic. CMV seropositive donors are still eligible to donate.
Leukocyte reduction is the removal of white blood cells by filtration. Leukoreduced blood products are less likely to cause HLA alloimmunization (development of antibodies against specific blood types), febrile non-hemolytic transfusion reaction, cytomegalovirus infection, and platelet-transfusion refractoriness.
Pathogen reduction treatment that involves, for example, the addition of riboflavin with subsequent exposure to UV light has been shown to be effective in inactivating pathogens (viruses, bacteria, parasites and white blood cells) in blood products. By inactivating white blood cells in donated blood products, riboflavin and UV light treatment can also replace gamma-irradiation as a method to prevent graft-versus-host disease (TA-GvHD).
Compatibility testing
Before a recipient receives a transfusion, compatibility testing between donor and recipient blood must be done. The first step before a transfusion is given is to type and screen the recipient's blood. Typing of recipient's blood determines the ABO and Rh status. The sample is then screened for any alloantibodies that may react with donor blood. It takes about 45 minutes to complete (depending on the method used). The blood bank scientist also checks for special requirements of the patient (e.g. need for washed, irradiated or CMV negative blood) and the history of the patient to see if they have previously identified antibodies and any other serological anomalies.
A positive screen warrants an antibody panel/investigation to determine if it is clinically significant. An antibody panel consists of commercially prepared group O red cell suspensions from donors that have been phenotyped for antigens that correspond to commonly encountered and clinically significant alloantibodies. Donor cells may have homozygous (e.g. K+k+), heterozygous (K+k-) expression or no expression of various antigens (K−k−). The phenotypes of all the donor cells being tested are shown in a chart. The patient's serum is tested against the various donor cells using an indirect Coombs test. Based on the reactions of the patient's serum against the donor cells, a pattern will emerge to confirm the presence of one or more antibodies. Not all antibodies are clinically significant (i.e. cause transfusion reactions, HDN, etc.). Once the patient has developed a clinically significant antibody it is vital that the patient receive antigen-negative red blood cells to prevent future transfusion reactions.
If there is no antibody present, an immediate spin crossmatch may be performed where the recipient serum and donor rbc are incubated. In the immediate spin method, two drops of patient serum are tested against a drop of 3–5% suspension of donor cells in a test tube and spun in a serofuge. Agglutination or hemolysis (i.e., positive Coombs test) in the test tube is a positive reaction. If the crossmatch is positive, then further investigation is needed. Patients with no history of red cell antibodies may qualify for computer-assisted crossmatch, which does not involve combining patient serum with donor cells.
If an antibody is suspected, potential donor units must first be screened for the corresponding antigen by phenotyping them. Antigen negative units are then tested against the patient plasma using an antiglobulin/indirect crossmatch technique at 37 degrees Celsius to enhance reactivity and make the test easier to read.
In urgent cases where crossmatching cannot be completed, and the risk of dropping hemoglobin outweighs the risk of transfusing uncrossmatched blood, O-negative blood is used, followed by crossmatch as soon as possible. O-negative is also used for children and women of childbearing age. It is preferable for the laboratory to obtain a pre-transfusion sample in these cases so a type and screen can be performed to determine the actual blood group of the patient and to check for alloantibodies.
Compatibility of ABO and Rh system for Red Cell (Erythrocyte) Transfusion
This chart shows possible matches in blood transfusion between donor and receiver using ABO and Rh system. The symbol indicates compatibility.
Adverse effects
In the same way that the safety of pharmaceutical products is overseen by pharmacovigilance, the safety of blood and blood products is overseen by haemovigilance. This is defined by the World Health Organization (WHO) as a system "...to identify and prevent occurrence or recurrence of transfusion related unwanted events, to increase the safety, efficacy and efficiency of blood transfusion, covering all activities of the transfusion chain from donor to recipient." The system should include monitoring, identification, reporting, investigation and analysis of adverse events, near-misses, and reactions related to transfusion and manufacturing. In the UK this data is collected by an independent organisation called SHOT (Serious Hazards Of Transfusion). Haemovigilance systems have been established in many countries with the objective of ensuring the safety of blood for transfusion, but their organisational set-up and operating principles can vary.
Transfusions of blood products are associated with several complications, many of which can be grouped as immunological or infectious. There is controversy on potential quality degradation during storage.
Immunologic reaction
Acute hemolytic reactions are defined according to Serious Hazards of Transfusion (SHOT) as "fever and other symptoms/signs of haemolysis within 24 hours of transfusion; confirmed by one or more of the following: a fall of Hb, rise in lactate dehydrogenase (LDH), positive direct antiglobulin test (DAT), positive crossmatch" This is due to destruction of donor red blood cells by preformed recipient antibodies. Most often this occurs because of clerical errors or improper ABO blood typing and crossmatching resulting in a mismatch in ABO blood type between the donor and the recipient. Symptoms include fever, chills, chest pain, back pain, hemorrhage, increased heart rate, shortness of breath, and rapid drop in blood pressure. When suspected, transfusion should be stopped immediately, and blood sent for tests to evaluate for presence of hemolysis. Treatment is supportive. Kidney injury may occur because of the effects of the hemolytic reaction (pigment nephropathy). The severity of the transfusion reaction is depended upon amount of donor's antigen transfused, nature of the donor's antigens, the nature and the amount of recipient antibodies.
Delayed hemolytic reactions occur more than 24 hours after a transfusion. They usually occur within 28 days of a transfusion. They can be due to either a low level of antibodies present prior to the start of the transfusion, which are not detectable on pre-transfusion testing; or development of a new antibody against an antigen in the transfused blood. Therefore, delayed haemolytic reaction does not manifest until after 24 hours when enough antibodies are available to cause a reaction. The red blood cells are removed by macrophages from the blood circulation into liver and spleen to be destroyed, which leads to extravascular haemolysis. This process usually mediated by anti-Rh and anti-Kidd antibodies. However, this type of transfusion reaction is less severe when compared to acute haemolytic transfusion reaction.
Febrile nonhemolytic reactions are, along with allergic transfusion reactions, the most common type of blood transfusion reaction and occur because of the release of inflammatory chemical signals released by white blood cells in stored donor blood or attack on donor's white blood cells by recipient's antibodies. This type of reaction occurs in about 7% of transfusions. Fever is generally short lived and is treated with antipyretics, and transfusions may be finished as long as an acute hemolytic reaction is excluded. This is a reason for the now-widespread use of leukoreduction — the filtration of donor white cells from red cell product units.
Allergic transfusion reactions are caused by IgE anti-allergen antibodies. When antibodies are bound to its antigens, histamine is released from mast cells and basophils. Either IgE antibodies from the donor's or recipient's side can cause the allergic reaction. It is more common in patients who have allergic conditions such as hay fever. Patient may feel itchy or having hives but the symptoms are usually mild and can be controlled by stopping the transfusion and giving antihistamines.
Anaphylactic reactions are rare life-threatening allergic conditions caused by IgA anti-plasma protein antibodies. For patients who have selective immunoglobulin A deficiency, the reaction is presumed to be caused by IgA antibodies in the donor's plasma. The patient may present with symptoms of fever, wheezing, coughing, shortness of breath, and circulatory shock. Urgent treatment with epinephrine is needed.
Post-transfusion purpura is an extremely rare complication that occurs after blood product transfusion and is associated with the presence of antibodies in the patient's blood directed against both the donor's and recipient's platelets HPA (human platelet antigen). Recipients who lack this protein develop sensitization to this protein from prior transfusions or previous pregnancies, can develop thrombocytopenia, bleeding into the skin, and can display purplish discolouration of skin which is known as purpura. Intravenous immunoglobulin (IVIG) is treatment of choice.
Transfusion-related acute lung injury (TRALI) is a syndrome that is similar to acute respiratory distress syndrome (ARDS), which develops during or within 6 hours of transfusion of a plasma-containing blood product. Fever, hypotension, shortness of breath, and tachycardia often occurs in this type of reaction. For a definitive diagnosis to be made, symptoms must occur within 6 hours of transfusion, hypoxemia must be present, there must be radiographic evidence of bilateral infiltrates and there must be no evidence of left atrial hypertension (fluid overload). It occurs in 15% of the transfused patient with mortality rate of 5 to 10%. Recipient risk factors includes: end-stage liver disease, sepsis, haematological malignancies, sepsis, and ventilated patients. Antibodies to human neutrophil antigens (HNA) and human leukocyte antigens (HLA) have been associated with this type of transfusion reaction. Donor's antibodies interacting with antigen positive recipient tissue result in release of inflammatory cytokines, resulting in pulmonary capillary leakage. The treatment is supportive.
Transfusion associated circulatory overload (TACO) is a common, yet underdiagnosed, reaction to blood product transfusion consisting of the new onset or exacerbation of three of the following within 6 hours of cessation of transfusion: acute respiratory distress, elevated brain natriuretic peptide (BNP), elevated central venous pressure (CVP), evidence of left heart failure, evidence of positive fluid balance, and/or radiographic evidence of pulmonary vascular congestion. Patients with congestive heart failure or kidney disease are more susceptible to volume overload. For especially vulnerable patients, a standard RBC unit could be split by sterile technique in the blood bank and administered over 8 hours instead of the standard 4 hours. Plasma transfusion is especially prone to causing TACO because large volumes are usually required to give any therapeutic benefit.
Transfusion-associated graft versus host disease frequently occurs in immunodeficient patients where recipient's body failed to eliminate donor's T cells. Instead, donor's T cells attack the recipient's cells. It occurs one week after transfusion. Fever, rash, diarrhoea are often associated with this type of transfusion reaction. Mortality rate is high, with 89.7% of the patients dead after 24 days. Immunosuppressive treatment is the most common way of treatment. Irradiation and leukoreduction of blood products is necessary for high risk patients to prevent T cells from attacking recipient cells.
Infection
The use of greater amount of red blood cells has been suggested to increase the risk of infections, not only transfusion-transmitted infections, but also due to a phenomenon known as transfusion-related immunomodulation (TRIM). TRIM may be caused by macrophages and their byproducts. In those who were given red blood cells only with significant anemia ("restrictive" strategy), serious infection rates were 10.6% while in those who were given red blood at milder levels of anemia ("liberal" strategy), serious infection rates were 12.7%.
On rare occasions, blood products are contaminated with bacteria. This can result in a life-threatening infection known as transfusion-transmitted bacterial infection. The risk of severe bacterial infection is estimated, , at about 1 in 2,500 platelet transfusions, and 1 in 2,000,000 red blood cell transfusions. Blood product contamination, while rare, is still more common than actual infection. The reason platelets are more often contaminated than other blood products is that they are stored at room temperature for short periods of time. Contamination is also more common with longer duration of storage, especially if that means more than 5 days. Sources of contaminants include the donor's blood, donor's skin, phlebotomist's skin, and containers. Contaminating organisms vary greatly, and include skin flora, gut flora, and environmental organisms. There are many strategies in place at blood donation centers and laboratories to reduce the risk of contamination. A definite diagnosis of transfusion-transmitted bacterial infection includes the identification of a positive culture in the recipient (without an alternative diagnosis) as well as the identification of the same organism in the donor blood.
Since the advent of HIV testing of donor blood in the mid/later 1980s, ex. 1985's ELISA, the transmission of HIV during transfusion has dropped dramatically. Prior testing of donor blood only included testing for antibodies to HIV. However, because of latent infection (the "window period" in which an individual is infectious, but has not had time to develop antibodies) many cases of HIV seropositive blood were missed. The development of a nucleic acid test for the HIV-1 RNA has dramatically lowered the rate of donor blood seropositivity to about 1 in 3 million units. As transmittance of HIV does not necessarily mean HIV infection, the latter could still occur at an even lower rate.
The transmission of hepatitis C via transfusion currently stands at a rate of about 1 in 2 million units. As with HIV, this low rate has been attributed to the ability to screen for both antibodies as well as viral RNA nucleic acid testing in donor blood.
Other rare transmissible infections include hepatitis B, syphilis, Chagas disease, cytomegalovirus infections (in immunocompromised recipients), HTLV, and Babesia.
Comparison table
Inefficacy
Transfusion inefficacy or insufficient efficacy of a given unit(s) of blood product, while not itself a "complication" per se, can nonetheless indirectly lead to complications – in addition to causing a transfusion to fully or partly fail to achieve its clinical purpose. This can be especially significant for certain patient groups such as critical-care or neonatals.
For red blood cells (RBC), by far the most commonly transfused product, poor transfusion efficacy can result from units damaged by the so-called storage lesion — a range of biochemical and biomechanical changes that occur during storage. With red cells, this can decrease viability and ability for tissue oxygenation. Although some of the biochemical changes are reversible after the blood is transfused, the biomechanical changes are less so, and rejuvenation products are not yet able to adequately reverse this phenomenon. There has been controversy about whether a given product unit's age is a factor in transfusion efficacy, specifically about whether "older" blood directly or indirectly increases risks of complications. Studies have not been consistent on answering this question, with some showing that older blood is indeed less effective but with others showing no such difference; these developments are being closely followed by hospital blood bankers — who are the physicians, typically pathologists, who collect and manage inventories of transfusable blood units.
Certain regulatory measures are in place to minimize RBC storage lesion — including a maximum shelf life (currently 42 days), a maximum auto-hemolysis threshold (currently 1% in the US, 0.8% in Europe), and a minimum level of post-transfusion RBC survival in vivo (currently 75% after 24 hours). However, all of these criteria are applied in a universal manner that does not account for differences among units of product. For example, testing for the post-transfusion RBC survival in vivo is done on a sample of healthy volunteers, and then compliance is presumed for all RBC units based on universal (GMP) processing standards (RBC survival by itself does not guarantee efficacy, but it is a necessary prerequisite for cell function, and hence serves as a regulatory proxy). Opinions vary as to the "best" way to determine transfusion efficacy in a patient in vivo. In general, there are not yet any in vitro tests to assess quality or predict efficacy for specific units of RBC blood product prior to their transfusion, though there is exploration of potentially relevant tests based on RBC membrane properties such as erythrocyte deformability and erythrocyte fragility (mechanical).
Physicians have adopted a so-called "restrictive protocol" — whereby transfusion is held to a minimum — in part because of the noted uncertainties surrounding storage lesion, in addition to the very high direct and indirect costs of transfusions. However, the restrictive protocol is not an option with some especially vulnerable patients who may require the best possible efforts to rapidly restore tissue oxygenation.
Although transfusions of platelets are far less numerous (relative to RBC), platelet storage lesion and resulting efficacy loss is also a concern.
Other
A relationship between intra-operative blood transfusion and cancer recurrence has been observed in colorectal cancer. In lung cancer intra-operative blood transfusion has been associated with earlier recurrence of cancer, worse survival rates and poorer outcomes after lung resection. Suppression of the immune system by blood transfusion has been implicated as playing a role in more than 10 different cancer types, through mechanisms involving the innate and adaptive immune system. Five major mechanisms for this include the lymphocyte-T set, myeloid-derived suppressor cells (MDSCs), tumor-associated macrophages (TAMs), natural killer cells (NKCs), and dendritic cells (DCs). Blood transfusion may modulate the activity of antitumor CD8+ cytotoxic T lymphocytes (CD8+/CTL), temporal response of Tregs, and the STAT3 signaling pathway. The role of the antitumor immune response in cancer therapeutics was explored historically through the use of bacteria to enhance the antitumor immune response and more recently in cellular Immunotherapy. However, the impact of transfusion-related immunomodulation (TRIM) on cancer progression has not been definitively established and requires further study.
In retrospective studies, blood transfusion has been associated with worse outcomes after cytoreductive surgery and HIPEC. However, correlation does not prove causation, and transfused patients often have more complicated surgeries and more underlying cardiopulmonary disease compared to untransfused patients; conclusions should be based on prospective randomized controlled trials.
Hypothermia can occur with transfusions with large quantities of blood products which normally are stored at cold temperatures. Core body temperature can go down as low as 32 °C and can produce physiologic disturbances. Prevention should be done with warming the blood to ambient temperature prior to transfusions. Blood warming devices are available to avoid the hemolysis that would occur from unsafe practices such as microwaving.
Transfusions with large amounts of red blood cells, whether due to severe hemorrhaging and/or transfusion inefficacy (see above), can lead to an inclination for bleeding. The mechanism is thought to be due to disseminated intravascular coagulation, along with dilution of recipient platelets and coagulation factors. Close monitoring and transfusions with platelets and plasma is indicated when necessary. Progressive hemorrhagic injury (PHI) in traumatic brain injured patients may be worsened by liberal transfusion strategies.
Metabolic alkalosis can occur with massive blood transfusions because of the breakdown of citrate stored in blood into bicarbonate. However, acidemia is common in massively transfused patients, and acid-base balance is affected by complex factors.
Hypocalcemia can also occur with massive blood transfusions because of the complex of citrate with serum calcium. Calcium levels below 0.9 mmol/L should be treated.
Blood doping has been used by athletes to increase physical stamina. A lack of knowledge and insufficient experience can turn a blood transfusion into a dangerous event. For example, improper storage involving freezing and thawing, or minor antigen incompatibility, could lead to hemolysis.
Frequency of use
Globally around 85 million units of red blood cells are transfused in a given year. The global demand is much higher and there is an unmet need for safe blood for transfusion in many low- and middle-income countries.
In the United States, blood transfusions were performed nearly 3 million times during hospitalizations in 2011, making it the most common procedure performed. The rate of hospitalizations with a blood transfusion nearly doubled from 1997, from a rate of 40 stays to 95 stays per 10,000 population. It was the most common procedure performed for patients 45 years of age and older in 2011, and among the top five most common for patients between the ages of 1 and 44 years.
According to the New York Times: "Changes in medicine have eliminated the need for millions of blood transfusions, which is good news for patients getting procedures like coronary bypasses and other procedures that once required a lot of blood." And, "Blood bank revenue is falling, and the decline may reach $1.5 billion a year this year [2014] from a high of $5 billion in 2008." In 2014, the Red Cross was predicting job losses as high as 12,000 within the next three to five years, roughly a quarter of the total in the industry. As of 2019, the trend of declining transfusions appeared to be stabilizing, with 10,852,000 RBC units transfused in the United States.
History
Beginning with William Harvey's experiments on the circulation of blood, recorded research into blood transfusion began in the 17th century, with successful experiments in transfusion between animals. However, successive attempts by physicians to transfuse animal blood into humans gave variable, often fatal, results.
Pope Innocent VIII is sometimes said to have been given "the world's first blood transfusion" by his Italian-Jewish physician Giacomo di San Genesio, who had him drink (by mouth) the blood of three 10-year-old boys. The boys consequently died, as did the Pope himself. However, the evidence for this story is unreliable and considered a possible anti-Jewish blood libel.
Early attempts
Animal blood
Working at the Royal Society in the 1660s, the physician Richard Lower began examining the effects of changes in blood volume on circulatory function and developed methods for cross-circulatory study in animals, obviating clotting by closed arteriovenous connections. The new instruments he was able to devise enabled him to perform the first reliably documented successful transfusion of blood in front of his distinguished colleagues from the Royal Society.
According to Lower's account, "...towards the end of February 1665 [I] selected one dog of medium size, opened its jugular vein, and drew off blood, until its strength was nearly gone. Then, to make up for the great loss of this dog by the blood of a second, I introduced blood from the cervical artery of a fairly large mastiff, which had been fastened alongside the first, until this latter animal showed ... it was overfilled ... by the inflowing blood." After he "sewed up the jugular veins", the animal recovered "with no sign of discomfort or of displeasure".
Lower had performed the first blood transfusion between animals. He was then "requested by the Honorable [Robert] Boyle ... to acquaint the Royal Society with the procedure for the whole experiment", which he did in December 1665 in the Society's Philosophical Transactions.
The first blood transfusion from animal to human was administered by Dr. Jean-Baptiste Denys, eminent physician to King Louis XIV of France, on June 15, 1667. He transfused the blood of a sheep into a 15-year-old boy, who survived the transfusion. Denys performed another transfusion into a labourer, who also survived. Both instances were likely due to the small amount of blood that was actually transfused into these people. This allowed them to withstand the allergic reaction.
Denys's third patient to undergo a blood transfusion was Swedish Baron Gustaf Bonde. He received two transfusions. After the second transfusion Bonde died. In the winter of 1667, Denys performed several transfusions on Antoine Mauroy with calf's blood. On the third account Mauroy died.
Six months later in London, Lower performed the first human transfusion of animal blood in Britain, where he "superintended the introduction in [a patient's] arm at various times of some ounces of sheep's blood at a meeting of the Royal Society, and without any inconvenience to him." The recipient was Arthur Coga, "the subject of a harmless form of insanity." Sheep's blood was used because of speculation about the value of blood exchange between species; it had been suggested that blood from a gentle lamb might quiet the tempestuous spirit of an agitated person and that the shy might be made outgoing by blood from more sociable creatures. Coga received 20 shillings () to participate in the experiment.
Lower went on to pioneer new devices for the precise control of blood flow and the transfusion of blood; his designs were substantially the same as modern syringes and catheters. Shortly after, Lower moved to London, where his growing practice soon led him to abandon research.
These early experiments with animal blood provoked a heated controversy in Britain and France. Finally, in 1668, the Royal Society and the French government both banned the procedure. The Vatican condemned these experiments in 1670. Blood transfusions fell into obscurity for the next 150 years.
Human blood
The science of blood transfusion dates to the first decade of the 20th century, with the discovery of distinct blood types leading to the practice of mixing some blood from the donor and the receiver before the transfusion (an early form of cross-matching).
In the early 19th century, British obstetrician Dr. James Blundell made efforts to treat hemorrhage by transfusion of human blood using a syringe. In 1818, after experiments with animals, he performed the first successful transfusion of human blood to treat postpartum hemorrhage. Blundell used the patient's husband as a donor, and extracted four ounces of blood from his arm to transfuse into his wife. During the years 1825 and 1830, Blundell performed 10 transfusions, five of which were beneficial, and published his results. He also invented a number of instruments for the transfusion of blood. He made a substantial amount of money from this endeavour, roughly $2 million ($50 million real dollars).
In 1840, at St George's Hospital Medical School in London, Samuel Armstrong Lane, aided by Blundell, performed the first successful whole blood transfusion to treat haemophilia.
However, early transfusions were risky and many resulted in the death of the patient. By the late 19th century, blood transfusion was regarded as a risky and dubious procedure, and was largely shunned by the medical establishment.
Work to emulate James Blundell continued in Edinburgh. In 1845 the Edinburgh Journal described the successful transfusion of blood to a woman with severe uterine bleeding. Subsequent transfusions were successful with patients of Professor James Young Simpson, after whom the Simpson Memorial Maternity Pavilion in Edinburgh was named.
Various isolated reports of successful transfusions emerged towards the end of the 19th century. The largest series of early successful transfusions took place at the Edinburgh Royal Infirmary between 1885 and 1892. Edinburgh later became the home of the first blood donation and blood transfusion services.
20th century
Only in 1901, when the Austrian Karl Landsteiner discovered three human blood groups (O, A, and B), did blood transfusion achieve a scientific basis and become safer.
Landsteiner discovered that adverse effects arise from mixing blood from two incompatible individuals. He found that mixing incompatible types triggers an immune response and the red blood-cells clump. The immunological reaction occurs when the receiver of a blood transfusion has antibodies against the donor blood-cells. The destruction of red blood cells releases free hemoglobin into the bloodstream, which can have fatal consequences. Landsteiner's work made it possible to determine blood group and allowed blood transfusions to take place much more safely. For his discovery he won the Nobel Prize in Physiology and Medicine in 1930; many other blood groups have been discovered since.
George Washington Crile is credited with performing the first surgery using a direct blood transfusion in 1906 at St. Alexis Hospital in Cleveland while a professor of surgery at Case Western Reserve University.
Jan Janský also discovered the human blood groups; in 1907 he classified blood into four groups: I, II, III, IV. His nomenclature is still used in Russia and in states of the former USSR, in which blood types O, A, B, and AB are respectively designated I, II, III, and IV.
Dr. William Lorenzo Moss's (1876–1957) Moss-blood typing technique of 1910 was widely used until World War II.
William Stewart Halsted, M.D. (1852–1922), an American surgeon, performed one of the first blood transfusions in the United States. He had been called to see his sister after she had given birth. He found her moribund from blood loss, and in a bold move withdrew his own blood, transfused his blood into his sister, and then operated on her to save her life.
Blood banks in WWI
While the first transfusions had to be made directly from donor to receiver before coagulation, it was discovered that by adding anticoagulant and refrigerating the blood it was possible to store it for some days, thus opening the way for the development of blood banks. John Braxton Hicks was the first to experiment with chemical methods to prevent the coagulation of blood at St Mary's Hospital, London in the late-19th century. His attempts, using phosphate of soda, however, proved unsuccessful.
The Belgian doctor Albert Hustin performed the first non-direct transfusion on March 27, 1914, though this involved a diluted solution of blood. The Argentine doctor Luis Agote used a much less diluted solution in November of the same year. Both used sodium citrate as an anticoagulant.
The First World War (1914–1918) acted as a catalyst for the rapid development of blood banks and transfusion techniques. Francis Peyton Rous and Joseph R. Turner at the Rockefeller University (then The Rockefeller Institute for Medical Research) made the first important discoveries that blood typing was necessary to avoid blood clumping (coagulation) and blood samples could be preserved using chemical treatment. Their first report in March 1915 showed that gelatine, agar, blood serum extracts, starch and beef albumin proved to be useless preservatives.
However, building on the same experiment, they discovered that a mixture sodium citrate and glucose (dextrose) solution was a perfect preservative; as they reported in February issue of the Journal of Experimental Medicine, the preserved bloods were just like fresh bloods and that they "function excellently when reintroduced into the body." Blood could be preserved for up to four weeks. An accompanying experiment using citrate-saccharose (sucrose) mixture was also a success which could maintain blood cells for two weeks. This use of citrate and sugars, sometimes known as Rous-Turner solution, was the foundation for the development of blood banks and improvement of transfusion method.
Another discovery of Rous and Turner was the most critical step in the safety of blood transfusion. Rous was well aware that Landsteiner's concept of blood types had not yet find practical value, as he remarked: "The fate of Landsteiner's effort to call attention to the practical bearing of the group differences in human bloods provides an exquisite instance of knowledge marking time on technique. Transfusion was still not done because (until at least 1915), the risk of clotting was too great." In June 1915, they made a crucial report in the Journal of the American Medical Association that agglutination could be avoided if the blood samples of the donor and recipient were tested before. Which they called a rapid and simple method for testing blood compatibility, sodium citrate was used to dilute the blood samples, and after mixing the recipient's and donor's blood in 9:1 and 1:1 parts, blood would either clump or remain watery after 15 minutes. According to their advice, blood without clumping "should always be chosen if possible."
Canadian doctor and Lieutenant Lawrence Bruce Robertson became instrumental in persuading the Royal Army Medical Corps to adopt the use of blood transfusion at the Casualty Clearing Stations for the wounded. In October 1915, Robertson performed his first wartime transfusion with a syringe to a patient who had multiple shrapnel wounds. He followed this up with four subsequent transfusions in the following months, and his success was reported to Sir Walter Morley Fletcher, director of the Medical Research Committee.
Robertson published his findings in the British Medical Journal in 1916 and, with the help of a few like-minded individuals (including the eminent physician Edward William Archibald), was able to persuade the British authorities of the merits of blood transfusion. Robertson went on to establish the first blood-transfusion apparatus at a Casualty Clearing Station on the Western Front in the spring of 1917. Robertson did not test crossmatching so that one died of hemolysis in his 1916 transfusion, and three in 1917.
Oswald Hope Robertson, a medical researcher and U.S. Army officer, was attached to the RAMC in 1917, where he became instrumental in establishing the first blood banks in preparation for the anticipated Third Battle of Ypres. He used sodium citrate as the anticoagulant; blood was extracted from punctures in the vein and was stored in bottles at British and American Casualty Clearing Stations along the Front. Robertson also experimented with preserving separated red blood cells in iced bottles. Geoffrey Keynes, a British surgeon, developed a portable machine that could store blood to enable transfusions to be carried out more easily.
Expansion
The secretary of the British Red Cross, Percy Lane Oliver, established the world's first blood-donor service in 1921. In that year, Oliver was contacted by King's College Hospital, where they were in urgent need of a blood donor.
After providing a donor, Oliver set about organizing a system for the voluntary registration of blood donors at clinics around London, with Sir Geoffrey Keynes appointed as a medical adviser. Volunteers were subjected to a series of physical tests to establish their blood group. The London Blood Transfusion Service was free of charge and expanded rapidly in its first few years of operation. By 1925 it was providing services for almost 500 patients; it was incorporated into the structure of the British Red Cross in 1926. Similar systems developed in other cities, including Sheffield, Manchester and Norwich, and the service's work began to attract international attention. France, Germany, Austria, Belgium, Australia and Japan established similar services.
Alexander Bogdanov founded an academic institution devoted to the science of blood transfusion in Moscow in 1925. Bogdanov was motivated, at least in part, by a search for eternal youth, and remarked with satisfaction on the improvement of his eyesight, suspension of balding, and other positive symptoms after receiving 11 transfusions of whole blood. Bogdanov died in 1928 as a result of one of his experiments, when the blood of a student with malaria and tuberculosis was given to him in a transfusion. Following Bogdanov's lead, Vladimir Shamov and Sergei Yudin in the USSR pioneered the transfusion of cadaveric blood from recently deceased donors. Yudin performed such a transfusion successfully for the first time on March 23, 1930, and reported his first seven clinical transfusions with cadaveric blood at the Fourth Congress of Ukrainian Surgeons at Kharkiv in September. However, this method was never used widely, even in the Soviet Union. Nevertheless, the Soviet Union was the first to establish a network of facilities to collect and store blood for use in transfusions at hospitals.
Frederic Durán-Jordà established one of the earliest blood banks during the Spanish Civil War in 1936. Duran joined the Transfusion Service at the Barcelona Hospital at the start of the conflict, but the hospital was soon overwhelmed by the demand for blood and the paucity of available donors. With support from the Department of Health of the Spanish Republican Army, Duran established a blood bank for the use of wounded soldiers and civilians. The 300–400 mL of extracted blood was mixed with 10% citrate solution in a modified Duran Erlenmeyer flask. The blood was stored in a sterile glass enclosed under pressure at 2 °C. During 30 months of work, the Transfusion Service of Barcelona registered almost 30,000 donors, and processed 9,000 liters of blood.
In 1937 Bernard Fantus, director of therapeutics at the Cook County Hospital in Chicago, established the first hospital blood-bank in the United States. In setting up a hospital laboratory that preserved, refrigerated and stored donor blood, Fantus originated the term "blood bank". Within a few years, hospital and community blood-banks were established across the United States. Until the middle of World War II, the newly established US blood banks rejected African-American donors. During the war, Black people were allowed to donate blood, but the donated blood was labeled as being suitable only for transfusion into another person from the same race.
Frederic Durán-Jordà fled to Britain in 1938 and worked with Dr Janet Vaughan at the Royal Postgraduate Medical School at Hammersmith Hospital to establish a system of national blood banks in London. With the outbreak of war appearing imminent in 1938, the War Office created the Army Blood Supply Depot (ABSD) in Bristol, headed by Lionel Whitby and in control of four large blood-depots around the country. British policy through the war was to supply military personnel with blood from centralized depots, in contrast to the approach taken by the Americans and Germans where troops at the front were bled to provide required blood. The British method proved more successful in adequately meeting all requirements, and over 700,000 donors were bled over the course of the war. This system evolved into the National Blood Transfusion Service established in 1946, the first national service to be implemented.
Stories tell of Nazis in Eastern Europe during World War II using captive children as repeated involuntary blood-donors.
Medical advances
A blood-collection program was initiated in the US in 1940 and Edwin Cohn pioneered the process of blood fractionation. He worked out the techniques for isolating the serum albumin fraction of blood plasma, which is essential for maintaining the osmotic pressure in the blood vessels, preventing their collapse.
Gordon R. Ward, writing in the correspondence columns of the British Medical Journal, proposed the use of blood plasma as a substitute for whole blood and for transfusion purposes as early as 1918. At the onset of World War II, liquid plasma was used in Britain. A large project, known as "Blood for Britain", began in August 1940 to collect blood in New York City hospitals for the export of plasma to Britain. A freeze-dried plasma package was developed by the Surgeons General of the Army and Navy, working with the National Research Council, which reduced breakage and made transportation, packaging, and storage much simpler.
The resulting dried plasma package came in two tin cans containing 400 mL bottles. One bottle contained enough distilled water to reconstitute the dried plasma contained within the other bottle. In about three minutes, the plasma would be ready to use and could stay fresh for around four hours. Dr. Charles R. Drew was appointed medical supervisor, and he was able to transform the test-tube methods into the first successful technique for mass production.
Another important breakthrough came in 1937–40 when Karl Landsteiner (1868–1943), Alex Wiener, Philip Levine, and R.E. Stetson discovered the Rhesus blood group system, which was found to be the cause of the majority of transfusion reactions up to that time. Three years later, the introduction by J.F. Loutit and Patrick L. Mollison of acid–citrate–dextrose (ACD) solution, which reduced the volume of anticoagulant, permitted transfusions of greater volumes of blood and allowed longer-term storage.
Carl Walter and W.P. Murphy Jr. introduced the plastic bag for blood collection in 1950. Replacing breakable glass bottles with durable plastic bags made from PVC allowed for the evolution of a collection system capable of safe and easy preparation of multiple blood components from a single unit of whole blood.
In the field of cancer surgery, the replacement of massive blood-loss became a major problem. The cardiac-arrest rate was high. In 1963 C. Paul Boyan and William S. Howland discovered that the temperature of the blood and the rate of infusion greatly affected survival rates, and introduced blood warming to surgery.
Further extending the shelf-life of stored blood up to 42 days was an anticoagulant preservative, CPDA-1, introduced in 1979, which increased the blood supply and facilitated resource-sharing among blood banks.
about 15 million units of blood products were transfused per year in the United States. By 2013 the number had declined to about 11 million units, because of the shift towards laparoscopic surgery and other surgical advances and studies that have shown that many transfusions were unnecessary. For example, the standard of care reduced the amount of blood transfused in one case from 750 to 200 mL. In 2019, 10,852,000 RBC units, 2,243,000 platelet units, and 2,285,000 plasma units were transfused in the United States.
Special populations
Neonate
To ensure the safety of blood transfusion to pediatric patients, hospitals are taking additional precautions to avoid infection and prefer to use pediatric blood units that are guaranteed "safe" from Cytomegalovirus. Some guidelines have recommended the provision of CMV-negative blood components and not simply leukoreduced components for newborns or low birthweight infants in whom the immune system is not fully developed, but practice varies. These requirements place additional restrictions on blood donors who can donate for neonatal use, which may be impractical given the rarity of CMV seronegative donors and the preference for fresh units.
Neonatal transfusions typically fall into one of two categories:
"Top-up" transfusions, to replace losses due to investigational losses and correction of anemia.
Exchange (or partial exchange) transfusions are done for removal of bilirubin, removal of antibodies and replacement of red cells (e.g., for anemia secondary to thalassemias and other hemoglobinopathies or fetal erythroblastosis).
Significant blood loss
A massive transfusion protocol is used when significant blood loss is present such as in major trauma, when more than ten units of blood are needed. Packed red blood cells, fresh frozen plasma, and platelets are generally administered. Typical ratios of fresh frozen plasma, platelets and packed red blood cells are between 1:1:1 and 1:1:2.
In some locations, blood has begun to be administered pre-hospital in an effort to reduce preventable deaths from significant blood loss. Earlier analyses suggested that in the US, up to 31,000 patients per year bleed to death that otherwise could have survived if pre-hospital transfusions were widely available. For example, when a mother experiences severe blood loss during pregnancy, ambulances are able to arrive with blood stored in portable, FDA listed blood refrigerators, similar to those found in blood banks. Once the infusion is given on scene, the patient and the ambulance have more time to get to a hospital for surgery and additional infusions if needed. This could be critical in rural areas or sprawling cities where patients can be far from a major hospital and the local emergency medical team may need to use blood infusions to keep that patient alive during transport. Larger studies pointed to improvements in 24-hour mortality with pre-hospital plasma and RBC transfusions, but no difference in 30-day or long-term mortality.
Unknown blood type
Because blood type O negative is compatible with anyone, it is often overused and in short supply. According to the Association for the Advancement of Blood and Biotherapies, the use of this blood should be restricted to persons with O negative blood, as nothing else is compatible with them, and women who might be pregnant and for whom it would be impossible to do blood group testing before giving them emergency treatment. Whenever possible, the AABB recommends that O negative blood be conserved by using blood type testing to identify a less scarce alternative.
Religious objections
Jehovah's Witnesses may object to blood transfusions because of their belief that blood is sacred.
Personal objections
Sometimes people refuse blood transfusions because of fears about the safety of the blood supply. Generally speaking, the rules of informed consent allow mentally competent adults to refuse blood transfusions even when their objections are based on misinformation or prejudice and even when their refusal may result in serious and permanent harm, including death. For example, since COVID-19 vaccines became available, some people in the US have refused blood transfusions because the donor might have been vaccinated, and they fear that this would cause secondhand harm to them. This choice is based on false beliefs, but a mentally competent adult's choices are normally respected. However, if doctor's opinions are that parents and guardians are making harmful choices about children, they can be overruled (in some jurisdictions) using legal arguments based on the harm principle; in this case if doctors believe that refusing the blood transfusion would put the child at risk of serious injury or death. Blood banks do not collect information that is irrelevant to the transfusion process, including the donors' race, ethnicity, sexual orientation, COVID-19 vaccination status, etc., so selecting blood units based on the individual's personal objections is not practical.
Research into alternatives
Although there are clinical situations where transfusion with red blood cells is the only clinically appropriate option, clinicians look at whether alternatives are feasible. This can be due to several reasons, such as patient safety, economic burden or scarcity of blood. Guidelines recommend blood transfusions should be reserved for patients with or at risk of cardiovascular instability due to the degree of their anaemia. In stable patients with iron deficiency anemia, oral or parenteral iron is recommended.
Thus far, there are no FDA-approved oxygen-carrying blood substitutes, which is the typical objective of a blood (RBC) transfusion. Non-blood volume expanders are available for cases where only volume restoration is required, but a substance with oxygen-carrying capacity would help doctors and surgeons avoid the risks of disease transmission and immune suppression, address the chronic blood donor shortage, and address the concerns of Jehovah's Witnesses and others who have religious objections to receiving transfused blood.
The research in this area is ongoing. A number of blood substitutes have been explored, but thus far they all have serious limitations. Most attempts to find a suitable alternative to blood thus far have concentrated on cell-free hemoglobin solutions. Blood substitutes could make transfusions more readily available in emergency medicine and in pre-hospital EMS care. If successful, such a blood substitute could save many lives, particularly in trauma where massive blood loss results. Hemopure, a hemoglobin-based therapy, is approved for use in South Africa and has been used in the United States on a case-by-case basis through the emergency Investigational New Drug (IND) process.
Veterinary use
Veterinarians also administer transfusions to other animals. Various species require different levels of testing to ensure a compatible match. For example, cats have 3 known blood types, cattle have 11, dogs have at least 13, pigs have 16, and horses over 30. However, in many species (especially horses and dogs), cross matching is not required before the first transfusion, as antibodies against non-self cell surface antigens are not expressed constitutively — i.e. the animal has to be sensitized before it will mount an immune response against the transfused blood.
The rare and experimental practice of inter-species blood transfusions (xenotransfusion) is a form of xenograft.
| Biology and health sciences | Medical procedures | null |
88873 | https://en.wikipedia.org/wiki/Nest | Nest | A nest is a structure built for certain animals to hold eggs or young. Although nests are most closely associated with birds, members of all classes of vertebrates and some invertebrates construct nests. They may be composed of organic material such as twigs, grass, and leaves, or may be a simple depression in the ground, or a hole in a rock, tree, or building. Human-made materials, such as string, plastic, cloth, or paper, may also be used. Nests can be found in all types of habitat.
Nest building is driven by a biological urge known as the nesting instinct in birds and mammals. Generally each species has a distinctive style of nest. Nest complexity is roughly correlated with the level of parental care by adults. Nest building is considered a key adaptive advantage among birds, and they exhibit the most variation in their nests ranging from simple holes in the ground to elaborate communal nests hosting hundreds of individuals. Nests of prairie dogs and several social insects can host millions of individuals.
Nest building
Purposes of nesting
Structural purposes
Nest building (nidification) is often driven by a biological urge in pregnant animals to protect one's offspring known as the nesting instinct. Animals build nests to protect their eggs, their offspring, or themselves from danger. The simplest nest structures are adapted to hide eggs from predators, shield them from the sun or other environmental factors, or simply keep them from being scattered in ocean currents. In some cases, nests also help provide safety in numbers for egg-laying animals.
Social purposes
Many nest builders provide parental care to their young, while others simply lay their eggs and leave. Brooding (incubating eggs by sitting on them) is common among birds. In general, nest complexity increases in relation to the level of parental care provided. Nest building reinforces social behavior, allowing for larger populations in small spaces to the point of increasing the carrying capacity of an environment. Insects that exhibit the most complex nidification also exhibit the greatest social structure. Among mammals, the naked mole-rat displays a caste structure similar to the social insects while building extensive burrows that house hundreds of individuals.
Usage of environment
Versatility in use of construction material may be an adaptive advantage (less energy used to gather materials) or a disadvantage (less ability to specialize construction). The available evidence suggests that natural selection more often favors specialization over flexibility in nest construction.
At the most basic level, there are only two types of nest building: sculpting and assembly.
Sculpting
Sculpting is the process of removing material to achieve the desired outcome. Most commonly this entails burrowing into the ground or plant matter to create a nesting site.
Assembly
Assembly entails gathering, transporting, and arranging materials to create a novel structure. Transportation has the greatest time and energy cost so animals are usually adapted to build with materials available in their immediate environment.
Building materials
Plant matter is the most common construction material for nests. Other common materials include fur or feathers, perhaps from the animal itself, mud or dirt, fecal matter, and specialized secretions from the animal's body.
Effects on environment
Nest building can have a substantial impact on the environment in which animals live. The combined digging activity of termites and mole-rats in South Africa has created a "mima prairie" landscape marked by huge areas of flat land punctuated by mounds wide and high. Similar structures exist in the United States, created by pocket gophers, and Argentina, rodents of the genus Ctenomys.
Lasting effects
Nests constructed by megapode birds have been mistaken for anthropological features by professionals, due to their exceptional height () and abundance (hundreds in a single location).
Nest builders
Nest architecture may be as useful for distinguishing species as the animals' physical appearance. Species identified through such means are called ethospecies. This is especially common in wasps and termites, but also can apply to birds. In most animals, there is some variation in nest construction between individuals. Whether these differences are driven by genetics or learned behavior is unknown.
With the exception of a few tunneling mammals, nest builders exhibit no specialized anatomy, instead making use of body parts primarily used for other purposes. This is possibly due to the sporadic nature of nest building, minimizing the selective pressures of anatomy used for nest building.
Birds
In general, birds are the most skilled nest builders, although not all species of birds build nests, some laying their eggs directly onto rock ledges or bare soil without first modifying the area. Complex nest building is considered to be one of the key adaptive advantages of birds. Nests help regulate temperature and reduce predation risks, thus increasing the chance that offspring live to adulthood.
Bird nests vary from simple depressions in the ground known as scrapes to largely unstructured collections of branches to elaborately woven pendants or spheres. The megapodes, one of the few groups who do not directly brood their young, incubate their young in a mound of decomposing vegetation. One species, Macrocephalon maleo, uses volcanic sand warmed by geothermal heat to keep its eggs warm. Among the simple nest builders are falcons, owls, and many shorebirds. The weavers exhibit perhaps the most elaborate nests, complete with strands of grass tied into knots. Most bird nests lie somewhere in the middle, with the majority building cup-shaped nests using some combination of mud, twigs and leaves, and feathers. Some birds, such flamingos and swifts, use saliva to help hold their nest together. The edible-nest swiftlet uses saliva alone to construct their nests. The rufous hornero nest is composed entirely of mud and feces, which is placed on tree branches to allow the sun to harden it into a usable structure. The tailorbirds stitch together leaves to provide cover for their nest sites.
The sociable weaver builds large communal nests in which many individual nests reside. They divide the nest using walls of grass placed atop a base of large sticks. At the entrances to the nest, sharp sticks are placed to ward off intruders. A single communal site can measure in height and in width. As many as 300 mating pairs may reside in the structure. Other birds often built their own nests on top of Weaver nest sites.
Some birds build nests in trees, some (such as eagles, vultures, and many seabirds) will build them on rocky ledges, and others nest on the ground or in burrows. Each species has a characteristic nest style, but few are particular about where they build their nests. Most species will choose whatever site in their environment best protects their nest, taking into account the nest's style. Several species will build on a cactus whenever possible. The bushtit and Bullock's oriole will suspend their nests from the tips of slender branches. The oropendolas take hanging nests to the extreme, constructing pouches up to tall using hanging vines as their base. The hanging nest is attached to thin tree branches, discouraging predation. Other species seek out crevices, using buildings or birdhouses when tree holes are not available.
Typical bird nests range from in size (hummingbirds) to (eagles) in diameter. The largest nest on record was made by a pair of bald eagles. It was in diameter, deep and was estimated to weigh more than . The lightest bird nests may weigh only a few grams. Incubation mounds of the mallee fowl can reach heights of and widths of . It is estimated the animal uses as much as of material in its construction. The extinct Sylviornis neocaledoniae may have constructed nesting mounds in diameter.
Mammals
Many species of small mammals such as rodents, rabbits, and ground squirrels dig burrows into the ground to protect themselves and their young. Prairie dogs build an elaborate system of tunnels which can span large stretches of land. One such structure, called a town, spanned and held an estimated 400 million individuals. Their homes are adapted to withstand large (above-ground) temperature variation, floods, and fire. Their young are raised in the deepest chambers where the temperature is the most stable.
Many mammals, including raccoons and skunks, seek natural cavities in the ground or in trees to build their nests. Raccoons, and some rodents, use leaves to build nests underground and in trees. Tree squirrels build their nests (dreys) in trees, while voles nest in tall grass. In some species, the nest serve as homes for adults while in others they are used to raise young. The duck-billed platypus and the echidna lay eggs in nests.
Gorillas build fresh nests daily out of leaves and other vegetation in which they sleep at night. They sometimes also build nests during the day for resting in. The smaller species of gorilla build their nests in trees, while the larger are confined to the ground. Nests of the western gorilla, the largest species, measure about in diameter.
Amphibians
Some species of frog build nests ranging from simple to modest complexity. Many stream-dwelling frogs lay their eggs in a gelatinous mass which they attach to underwater vegetation to prevent eggs from washing away. Nests can have other protective qualities. For example, the female Fletcher frog beats secreted mucus into a froth, creating a structure that serves as a line of defense against thermal extremes, predation, and desiccation.
Fish
Fish engage in nest building activities ranging from simply scooping out sediment to building enclosed structures out of plant matter. Male sticklebacks produce a special enzyme in their kidneys that they use to bind plants together.
Reptiles
The American alligator is known for its parenting skills. They build large nests of mud and vegetation on river banks or vegetation mats. The female digs a hole in the center to lay her eggs, covers them, and then guards them for two months until they hatch. When eggs start to hatch, she breaks open the nest which has hardened over time and leads the young to the water where she continues to care for them for another year. Alligators are very particular about their nesting sites and will abandon a site if things go wrong.
Cobras use leaves and other debris to build nests in which they lay eggs that both sexes guard. They carry the vegetation to the nest site by kinking their necks. Sea turtles dig a hole in the sand above the high tide line in which they lay their eggs. They then cover the soft eggs to protect them from the sun and predators and leave.
Dinosaurs
From the fossil record, it is known that many, or perhaps all, dinosaurs laid eggs. Paleontologists have identified a number of features that allow them to distinguish a nesting site from a random clustering of eggs. Those include regular clustering patterns, the co-occurrence of whole eggs with broken eggs and/or hatchlings, and the occurrence of physical features such as evidence of excavation.
The Oviraptor nests of Mongolia are perhaps the most famous case of dinosaur nesting. One specimen was found fossilized atop a nest in a brooding posture, proving the animal had been poorly named (Oviraptor means "egg taker").
A site known as Egg Mountain in Montana provides exceptional evidence of dinosaur nesting behavior. The site features dozens of nests each with 20 or more eggs belonging to the Maiasaura. Juvenile teeth at the site exhibit signs of wear, while the leg bones are not developed enough to walk. This allowed scientists to conclude that the species provided extensive parental care for its young. It is likely the species covered its nests with sand and vegetation to keep them warm and nested in colonies for increased protection.
Insects
Social insects, including most species of ants, bees, termites, and wasps, are nest builders. Their often elaborate nests may be found above or below ground. Features often include ventilation systems and separate chambers for the queen, her eggs, and developing individuals.
Bees and hornets often seek out natural cavities in which to construct their nests, also known as hives, in which they store food and raise their young. Other species of bee and some wasps dig holes in the ground or chew through wood. In the species Megachile rotundata, for example, females construct tubular-shaped nests in rotting wood as well as small holes in the ground, creating, each cell made from circular disks cut from plant leaves using the bee's mandibles. Bee nests are founded upon the wax the secrete from their bodies, while those of wasps are dependent on their ability to turn plant water into paper using their saliva. Nests often exhibit divided living, with eggs and food stores kept in distinct parts of the hive. Vespid wasps build complex nests from paper-like material where they lay eggs in individual cells. When the young hatch, their parents feed them chewed up larvae. Different species exhibit different nest structures. Paper wasp nesting consist of a single tier of cells, while yellow jacket nests can be many layers thick, reaching up to in diameter. Nesting strategies can be plastic, for instance the wasp Parischnogaster mellyi will significant vary its nest construction based on environmental conditions, and the wasp Mischocyttarus mexicanus is known to nest in groups or alone depending on the distribution of potential nest sites in the area. Nest sizes vary dramatically and the largest wasp nest on record measured in diameter and was tall. Found in New Zealand, it was likely built by the German wasp.
Termites build elaborate nests that span multiple generations and may last decades. Using chewed wood, mud, and feces they build large mounds which may extend well into the air. The largest nests, built by members of genus Amitermes, stand nearly tall with a similar circumference at the base, and host millions of individuals. Termite mounds are constructed to allow for excellent air flow, regulating the mound temperature. The mounds protect against drying and predation allowing many species to lose ancestral traits such as hard bodies, skin pigmentation, and good eyesight. Magnetic termites construct their nests with flattened sides along the north–south axis to ensure maximum warming during the winter, while exposing minimal surface area to the harshest mid-day sunshine. Other termite species use their nests to farm fungi.
Ant nests feature an elaborate colony structure that may extend or more underground. As the structure gets further underground, individual chambers become farther and farther apart indicating that the ant is aware of its depth. It is hypothesized that they accomplish this by sensing the level of carbon dioxide in the soil. The leaf cutter ant builds a complex nest which can house 8 million individuals. Its nests feature numerous chambers, most notably garden chambers where they farm fungus on leaves they harvest from the forest.
Species such as the carpenter ant and the wasp Polistes exclamans build "satellite nests" - smaller nests near, but separate from, the main nest. These satellite nests are used as an insurance against predators and parasites; if the original nest is attacked, surviving members can move the satellite nest. Other species such as the Black hover wasp, Parischnogaster alternata, construct nests in clusters with the central core composed of older colonies surrounded by younger colonies.
The Eastern carpenter bee, Xylocopa virginica, is unique in that individuals of that species build their nests in wood, bamboo culms, agave stalks, and other similar materials, although their preferred nesting material is pine or cedar lumber. When digging the nests, they use the wood shavings scraped from the wall to create partitions within the tunnels. The nests are usually round and have about 1-4 tunnels, each with multiple branches. Because these materials are often useful for humans in construction, X. virgininica's nesting behavior presents the disadvantage of weakening wood in manmade structures.
Effects on other species
The abundance of biological resources within the nest has led to a number of specialized predators. The aardvark and the ant eater use long tongues to prey upon termite and ant nests. Birds such as the honey buzzard specialize on wasp and bee nests, a resource also targeted by the tropical hornet. Symbiosis, ranging from feeding on waste to obligate parasitism, is common within the nest. Ant nests alone support symbiotes spanning six classes of arthropods which includes 35 families just from the beetles.
Names of nests
A badger's nest is called a sett.
A beaver's nest is called a lodge.
An eagle's nest is called an eyrie.
An otter's den is called a holt or a couch.
A pheasant's nest is called a nide.
A rabbit's nest is called a form.
A squirrel's or ringtail possum's nest is called a drey.
A wasp's nest is called a vespiary.
| Biology and health sciences | Shelters and structures | Animals |
88891 | https://en.wikipedia.org/wiki/Tilapia | Tilapia | Tilapia ( ) is the common name for nearly a hundred species of cichlid fish from the coelotilapine, coptodonine, heterotilapine, oreochromine, pelmatolapiine, and tilapiine tribes (formerly all were "Tilapiini"), with the economically most important species placed in the Coptodonini and Oreochromini. Tilapia are mainly freshwater fish inhabiting shallow streams, ponds, rivers, and lakes, and less commonly found living in brackish water. Historically, they have been of major importance in artisanal fishing in Africa, and they are of increasing importance in aquaculture and aquaponics. Tilapia can become a problematic invasive species in new warm-water habitats such as Australia, whether deliberately or accidentally introduced, but generally not in temperate climates due to their inability to survive in cold water.
Traditionally a popular and affordable food in the Philippines with a mild taste, tilapia has been the fourth-most consumed fish in the United States since 2002, favored for its low cost and easy preparation. It is commonly fried or broiled as part of a dish.
Etymology
The common name "tilapia" is based on the name of the cichlid genus Tilapia, which is itself a latinization of either tlhapi, the Tswana word for "fish", or the Greek word "tilon," referring to a fish mentioned by Aristotle, combined with "apios," meaning "distant". Scottish zoologist Andrew Smith named the genus in 1840.
History
The aquaculture of Nile tilapia dates from Ancient Egypt, where it was represented by the hieroglyph K1, of the Gardiner list: 𓆛.
The fish symbolized rebirth in Egyptian art, and was associated in Egyptian belief-systems with Hathor, goddess of fertility and abundance.
The tilapia was also said to accompany and protect the sun god on his daily journey across the sky, and a tilapia painted on tomb walls may recall part of spell 15 of the Book of the Dead by which the deceased hopes to take his place in the sun boat: "You see the tilapia in its [true] form at the turquoise pool", and "I behold the tilapia in its [true] nature guiding the speedy boat in its waters."
Tilapia were one of the three main types of fish caught in Talmudic times from the Sea of Galilee, specifically the Galilean comb (Sarotherodon galilaeus). Today, in Modern Hebrew, the fish species is called amnoon (probably a compound of am, "mother" and noon, "fish"). In English, it is sometimes known by the name "St. Peter's fish", which comes from the narrative in the Gospel of Matthew about the apostle Peter catching a fish that carried a coin in its mouth. Though the passage does not name the fish, different tilapia species (Sarotherodon galilaeus, Oreochromis aureus, Coptodon zillii, and Tristramella) are found in the Sea of Galilee, where the author of the Gospel of Matthew recounts the event took place. These species have been the target of small-scale artisanal fisheries in the area for thousands of years.
Characteristics
Tilapia typically have laterally compressed, deep bodies. Like other cichlids, their lower pharyngeal bones are fused into a single tooth-bearing structure. A complex set of muscles allows the upper and lower pharyngeal bones to be used as a second set of jaws for processing food (cf. morays), allowing a division of labor between the "true jaws" (mandibles) and the "pharyngeal jaws". This means they are efficient feeders that can capture and process a wide variety of food items. Their mouths are protrusible, usually bordered with wide and often swollen lips. The jaws have conical teeth. Typically, tilapia have a long dorsal fin, and a lateral line that often breaks towards the end of the dorsal fin, and starts again two or three rows of scales below. Some Nile tilapia can grow as long as .
Other than their temperature sensitivity, tilapia exist in or can adapt to a very wide range of conditions. An extreme example is the Salton Sea, where tilapia introduced when the water was merely brackish now live in salt concentrations so high that other marine fish cannot survive.
Tilapia are also known to be mouth-brooding species, which means they carry the fertilized eggs and young fish in their mouths for several days after the yolk sac is absorbed.
Gallery
Species
Historically, all tilapia have been included in their namesake genus Tilapia. In recent decades, some were moved into a few other genera, notably Oreochromis, and Sarotherodon. Even with this modification, apparently Tilapia was strongly poly– or paraphyletic. In 2013, a major taxonomic review resolved this by moving most former Tilapia spp. to several other genera. As a consequence, none of the species that are of major economic importance remain in Tilapia, but are instead placed in Coptodon, Oreochormis, and Sarotherodon.
Exotic and invasive species
Tilapia have been used as biological controls for certain aquatic plant problems. They have a preference for a floating aquatic plant, duckweed (Lemna spp.), but also consume some filamentous algae. In Kenya, tilapia were introduced to control mosquitoes, which were causing malaria, because they consume mosquito larvae, consequently reducing the numbers of adult female mosquitoes, the vector of the disease. These benefits are, however, frequently outweighed by the negative aspects of tilapia as invasive species.
Tilapia are unable to survive in temperate climates because they require warm water. The pure strain of the blue tilapia, Oreochromis aureus, has the greatest cold tolerance and dies at , while all other species of tilapia die at a range of . As a result, they cannot invade temperate habitats and disrupt native ecologies in temperate zones; however, they have spread widely beyond their points of introduction in many fresh and brackish tropical and subtropical habitats, often disrupting native species significantly. Because of this, tilapia are on the IUCN's 100 of the World's Worst Alien Invasive Species list. In the United States, tilapia are found in much of the south, especially Florida and Texas, and as far north as Idaho, where they survive in power-plant discharge zones. Tilapia are also currently stocked in the Phoenix, Arizona, canal system as an algal growth-control measure. In a Washington, D.C. fishing report from 21 June 2024, it was reported that an angler caught a Tilapia on a crankbait at the Jones Point Park under the Woodrow Wilson Bridge, which is on the Potomac River. Many state fish and wildlife agencies in the United States, Australia, South Africa, and elsewhere consider them to be invasive species.
Aquarium species
Larger tilapia species are generally poor community aquarium fish because they eat plants, dig up the bottom, and fight with other fish. The larger species are often raised as a food source, though, because they grow rapidly and tolerate high stocking densities and poor water quality.
Smaller West African species, such as Coelotilapia joka and species from the crater lakes of Cameroon, are more popular as aquarium fish. In specialised cichlid aquaria, tilapia can be mixed successfully with nonterritorial cichlids, armored catfish, tinfoil barbs, garpike, and other robust fish. Some species, including Heterotilapia buttikoferi, Coptodon rendalli, Pelmatolapia mariae, C. joka, and the brackish-water Sarotherodon melanotheron, have attractive patterns and are quite decorative.
Commercial species
Tilapia were originally farmed in their native Africa and Levant. Fast-growing, tolerant of stocking density, and adaptable, tilapia have been introduced to and are farmed extensively in many parts of Asia and are increasingly common aquaculture targets elsewhere.
Aquaculture
Farmed tilapia production in 2002 worldwide was about annually, with an estimated value of US$1.8 billion, about equal to those of salmon and trout.
Unlike carnivorous fish, tilapia can feed on algae or any plant-based food. This reduces the cost of tilapia farming, reduces fishing pressure on prey species, avoids concentrating toxins that accumulate at higher levels of the food chain, and makes tilapia the preferred "aquatic chickens" of the trade.
Because of their large size, rapid growth, and palatability, tilapia cichlids are the focus of major farming efforts, specifically various species of Oreochromis, Sarotherodon, and Coptodon (all were formerly in the namesake genus Tilapia). Like other large fish, they are a good source of protein and popular among artisanal and commercial fisheries. Most such fisheries were originally found in Africa, but outdoor fish farms in tropical countries, such as Papua New Guinea, the Philippines, and Indonesia, are underway in freshwater lakes. In temperate zone localities, tilapiine farming operations require energy to warm the water to tropical temperatures. One method uses waste heat from factories and power stations.
At 1.3 million tonnes per annum, China is the largest tilapia producer in the world, followed by Egypt with 0.5 million. The US, by comparison, produces 10 thousand tonnes against a consumption of 2.5 million.
In modern aquaculture, wild-type Nile tilapia are not too often seen, as the dark color of their flesh is not much desired by many customers, and because it has a bit of a reputation of being a rough fish associated with poverty. However, they are fast-growing and give good fillets; leucistic ("red") breeds which have lighter meat have been developed and are very popular.
Hybrid stock is also used in aquaculture; Nile × blue tilapia hybrids are usually rather dark, but a light-colored hybrid breed known as "Rocky Mountain White" tilapia is often grown due to its very light flesh and tolerance of low temperatures.
Commercially grown tilapia are almost exclusively male, typically done by adding male sex hormone in the food to the tilapia fry, causing any potential female tilapia to change sex to male. It can also be achieved through hybridization of certain tilapia species or the use of so-called "supermales" that have homozygous male sex chromosomes (resulting in all their offspring receiving a male sex chromosome and thus becoming males). Males are preferred because they grow much faster than females. Additionally, because tilapia are prolific breeders, the presence of female tilapia results in rapidly increasing populations of small fish, rather than a stable population of harvest-size animals.
Tilapia, as a "traditional" dish, was first introduced in the Philippines in 1950 with Mozambique tilapia followed by Nile tilapia in 1972 from Thailand. Strains such as terapon (gunggong) and white goby (biyang puti) were abundant in Laguna de Bay. In 1988, WorldFish, Filipino and Norwegian researchers started aquaculture of the Genetically Improved Farmed Tilapia (GIFT).
Other methods of tilapia population control are polyculture, with predators farmed alongside tilapia or hybridization with other species.
As food
Tilapia, introduced to the state of Tamil Nadu, India in the 1950s as a cheap protein source, quickly spread across the state's freshwater bodies because of its ability to thrive in polluted, low-oxygen environments.
The Fisheries Research Committee has recommended the culture of this fish only in areas in the western slopes of the Western Ghats and the coastal strip between Cape Comorin and south of and excluding river Tapti, and in districts of Tinnevellv, Madurai and Ramnad in Madras south of and including river Vaigai. The committee has suggested further investigation to enable it to examine the question of desirability of Tilapia culture in other parts of India.
Whole tilapia fish can be processed into skinless, boneless fillets. In some of the commercial strains, the yield has been reported up to 47% at harvest weight.
Tilapia are among the commercially important aquaculture species that are susceptible to off-flavors (others include trout, barramundi, and channel catfish). These 'muddy' or 'musty' flavors are normally caused by geosmin and 2-methylisoborneol, organic products of ubiquitous cyanobacteria that are often present or bloom sporadically in water bodies and soil. These flavors are no indication of freshness or safety of the fish, but they make the product unattractive to consumers. Simple quality-control procedures are known to be effective in ensuring the quality of fish entering the market.
Tilapia have very low levels of mercury. Tilapia are low in saturated fat, calories, carbohydrates, and sodium, and are a good protein source. They also contain the micronutrients phosphorus, niacin, selenium, vitamin B12, and potassium.
Tilapia may be a less nutritious fish than generally believed. The fish's omega-3 fatty acid content is often far lower than that of other commonly eaten fish species. Their omega-6 fatty acid levels are unusually high. Multiple studies have evaluated the effects of adding flaxseed derivatives (a vegetable source of omega-3 fatty acids) to the feed of farmed tilapia. These studies have found both the more common omega-3 fatty acid found in the flax, ALA and the two types almost unique to animal sources (DHA and EPA), increased in the fish fed this diet. Guided by these findings, tilapia farming techniques could be adjusted to address the nutritional criticisms directed at the fish, while retaining its advantage as an omnivore capable of feeding on economically and environmentally inexpensive vegetable protein. Adequate diets for salmon and other carnivorous fish can alternatively be formulated from protein sources such as soybean, although soy-based diets with soy oil may also change in the balance between omega-6 and omega-3 fatty acids.
Miscellaneous uses
Ecological agent
Tilapia serve as a natural, biological control for most aquatic plant problems. They consume floating aquatic plants, such as duckweed watermeal (Lemna spp.), most "undesirable" submerged plants, and most forms of algae. In the United States and countries such as Thailand, they are becoming the plant-control method of choice, reducing or eliminating the use of toxic chemicals and heavy metal-based algaecides. However, their environmental impact as an invasive species may outweigh their ecological benefit. Their tolerance for a wide range of environmental conditions allow them to thrive in polluted or otherwise degraded aquatic habitats.
They are used for zoo ponds as a source of food for birds.
Tilapia can be farmed with shrimp in a symbiotic manner, positively enhancing the productive output of both.
Arkansas stocks many public ponds and lakes to help with vegetation control, favoring tilapia as a robust forage species and for anglers.
In Kenya, tilapia help control mosquitoes, which carry malaria parasites. They consume mosquito larvae, which reduces the numbers of adult females, the disease's vector.
Medical use
In Brazil, Nile tilapia (Oreochromis niloticus) fish skin applied as a bandage is first used in a 2017 clinical trial to treat burn injuries, after successful trial in rats. In the United States, tilapia skin has been used to successfully treat third-degree wounds to the paws of two black bears caught in California's Thomas wildfire, and also to treat burns on the paws of a black bear from California's Carr wildfire. Nile tilapia skin has completed a phase III clinical trial for superficial partial-thickness burns. The fish skin group showed faster healing, lower pain, reduced dressing changes, and lower treatment costs compared to silver sulfadiazine cream control.
Nile tilapia skin has also been used in neovaginoplasty as a skin graft material, for Müllerian agenesis, vaginal stenosis, and gender-affirming surgery.
The skin-growing properties of tilapia skin are believed to be linked to its high type I collagen content and structural similarities to human skin. The material is also quite strong despite its low thickness. The current procedure for skin use calls for chemical sterilization with chlorhexidine, immersion in glycerol, followed by gamma ray sterilization, plus a few washes with saline in between.
Parasites
As with most fish, tilapia harbor a variety of parasites. For the monogeneans, these especially include species of the megadiverse genus Cichlidogyrus, which are gill parasites. Species of Enterogyrus are parasites in the digestive system. Tilapia, as important aquaculture fishes, have been introduced widely all over the world, and often carried their monogenean parasites with them. In South China, a 2019 study has shown that nine species of monogeneans were carried by introduced tilapia.
| Biology and health sciences | Acanthomorpha | null |
88931 | https://en.wikipedia.org/wiki/Cataract | Cataract | A cataract is a cloudy area in the lens of the eye that leads to a decrease in vision of the eye. Cataracts often develop slowly and can affect one or both eyes. Symptoms may include faded colours, blurry or double vision, halos around light, trouble with bright lights, and difficulty seeing at night. This may result in trouble driving, reading, or recognizing faces. Poor vision caused by cataracts may also result in an increased risk of falling and depression. Cataracts cause 51% of all cases of blindness and 33% of visual impairment worldwide.
Cataracts are most commonly due to aging but may also occur due to trauma or radiation exposure, be present from birth, or occur following eye surgery for other problems. Risk factors include diabetes, longstanding use of corticosteroid medication, smoking tobacco, prolonged exposure to sunlight, and alcohol. The underlying mechanism involves accumulation of clumps of protein or yellow-brown pigment in the lens that reduces transmission of light to the retina at the back of the eye. Diagnosis is by an eye examination.
Wearing sunglasses and a wide brimmed hat, eating leafy vegetables and fruits, and avoiding smoking may reduce the risk of developing cataracts, or slow the process. Early on, the symptoms may be improved with glasses. If this does not help, surgery to remove the cloudy lens and replace it with an artificial lens is the only effective treatment. Cataract surgery is not readily available in many countries, and surgery is needed only if the cataracts are causing problems and generally results in an improved quality of life.
About 20 million people worldwide are blind due to cataracts. It is the cause of approximately 5% of blindness in the United States and nearly 60% of blindness in parts of Africa and South America. Blindness from cataracts occurs in about 10 to 40 per 100,000 children in the developing world, and 1 to 4 per 100,000 children in the developed world. Cataracts become more common with age. In the United States, cataracts occur in 68% of those over the age of 80 years. Additionally they are more common in women, and less common in Hispanic and Black people.
Signs and symptoms
Signs and symptoms vary depending on the type of cataract, though considerable overlap occurs. People with nuclear sclerotic or brunescent cataracts often notice a reduction of vision. Nuclear cataracts typically cause greater impairment of distance vision than of near vision. Those with posterior subcapsular cataracts usually complain of glare as their major symptom.
The severity of cataract formation, assuming no other eye disease is present, is judged primarily by a visual acuity test. Other symptoms include frequent changes of glasses and colored halos due to hydration of lens.
Congenital cataracts can result in amblyopia if not treated in a timely manner.
Causes
Age
Age is the most common cause of cataracts. Lens proteins denature and degrade over time, and this process is accelerated by diseases such as diabetes mellitus and hypertension. Environmental factors, including toxins, radiation, and ultraviolet light have cumulative effects which are worsened by the loss of protective and restorative mechanisms due to alterations in gene expression and chemical processes within the eye.
Oxidative stress associated with lipid peroxidation is an important pathogenic mechanism in cataract formation. Senile cataracts are associated with a decrease in antioxidant capacity in the lens. An increase in oxidative stress in the lens or a decrease in the ability to remove reactive oxygen species can lead to the lens becoming more opaque.
Trauma
Blunt trauma causes swelling, thickening, and whitening of the lens fibers. While the swelling normally resolves with time, the white color may remain. In severe blunt trauma, or in injuries that penetrate the eye, the capsule in which the lens sits can be damaged. This damage allows fluid from other parts of the eye to rapidly enter the lens leading to swelling and then whitening, obstructing light from reaching the retina at the back of the eye. Cataracts may develop in 0.7 to 8.0% of cases following electrical injuries. Blunt trauma can also result in star- (stellate) or petal-shaped cataracts.
Radiation
Cataracts can arise as an effect of exposure to various types of radiation. X-rays, one form of ionizing radiation, may damage the DNA of lens cells. Ultraviolet light, specifically UVB, has also been shown to cause cataracts, and some evidence indicates sunglasses worn at an early age can slow its development in later life. Microwaves, a type of nonionizing radiation, may cause harm by denaturing protective enzymes (e.g., glutathione peroxidase), by oxidizing protein thiol groups (causing protein aggregation), or by damaging lens cells via thermoelastic expansion. The protein coagulation caused by electric and heat injuries whitens the lens. This same process is what makes the clear albumen of an egg become white and opaque during cooking.
Genetics
The genetic component is strong in the development of cataracts, most commonly through mechanisms that protect and maintain the lens. The presence of cataracts in childhood or early life can occasionally be due to a particular syndrome. Examples of chromosome abnormalities associated with cataracts include 1q21.1 deletion syndrome, cri-du-chat syndrome, Down syndrome, Patau's syndrome, trisomy 18 (Edward's syndrome), and Turner's syndrome, and in the case of neurofibromatosis type 2, juvenile cataract on one or both sides may be noted. Examples of single-gene disorder include Alport's syndrome, Conradi's syndrome, cerebrotendineous xanthomatosis, myotonic dystrophy, and oculocerebrorenal syndrome or Lowe syndrome.
Skin diseases
The skin and the lens have the same embryological origin and so can be affected by similar diseases. Those with atopic dermatitis and eczema occasionally develop shield ulcer cataracts. Ichthyosis is an autosomal recessive disorder associated with cuneiform cataracts and nuclear sclerosis. Basal-cell nevus and pemphigus have similar associations.
Smoking and alcohol
Cigarette smoking has been shown to increase the risk of age-related cataract and nuclear cataract. Evidence is conflicting over the effect of alcohol. Some surveys have shown a link, but others that followed people over longer terms have not.
Inadequate vitamin C
Low vitamin C intake and serum levels have been associated with greater cataract rates. However, use of supplements of vitamin C has not demonstrated benefit.
Medications
Some medications, such as systemic, topical, or inhaled corticosteroids, may increase the risk of cataract development. Corticosteroids most commonly cause posterior subcapsular cataracts. People with schizophrenia often have risk factors for lens opacities (such as diabetes, hypertension, and poor nutrition). Second-generation antipsychotic medications are unlikely to contribute to cataract formation. Miotics and triparanol may increase the risk.
Post-operative
Nearly every person who undergoes a vitrectomy—without ever having had cataract surgery—will experience progression of nuclear sclerosis after the operation. This may be because the native vitreous humor is different from the solutions used to replace the vitreous (vitreous substitutes), such as BSS Plus. This may also be because the native vitreous humour contains ascorbic acid which helps neutralize oxidative damage to the lens and because conventional vitreous substitutes do not contain ascorbic acid. Accordingly, for phakic patients requiring a vitrectomy it is becoming increasingly common for ophthalmologists to offer the vitrectomy combined with prophylactic cataract surgery to prevent cataract formation.
Hyperbaric oxygen therapy
Hyperbaric oxygen therapy (HBOT) is the administration of 100% oxygen at pressures greater than one-atmosphere absolute pressure (1 ATA) for a therapeutic purpose. HBOT can have several side effects, including the long-term development of cataracts. This is rare and generally associated with multiple HBOT exposures over a long period. As it does not usually become symptomatic during HBOT, it may often go unrecognised and is probably under-reported. Evidence is emerging that lifetime dosage of oxygen may be a precipitating factor in the development of age-related cataracts. Nuclear cataracts have been hypothesized to be the end stage of the far better known phenomenon of hyperbaric myopic shift.
Other diseases
Metabolic and nutritional diseases
Aminoaciduria or Lowe's syndrome
Cerebrotendineous xanthomatosis
Diabetes mellitus
Fabry's disease
Galactosemia / galactosemic cataract
Homocystinuria
Hyperparathyroidism
Hypoparathyroidism
Hypervitaminosis D
Hypothyroidism
Hypocalcaemia
Mucopolysaccharidoses
Wilson's disease
Congenital
Congenital syphilis
Cytomegalic inclusion disease
Rubella
Cockayne syndrome
Genetic syndromes
Down syndrome
Patau syndrome
Edwards syndrome
Infections:
Cysticercosis
Leprosy
Onchocerciasis
Toxoplasmosis
Varicella
Secondary to other eye diseases:
Retinopathy of prematurity
Aniridia
Uveitis
Retinal detachment
Retinitis pigmentosa
Diagnosis
Classification
Cataracts may be partial or complete, stationary or progressive, hard or soft. Histologically, the main types of age-related cataracts are nuclear sclerosis, cortical, and posterior subcapsular.
is the most common type of cataract, and involves the central or 'nuclear' part of the lens. This eventually becomes hard, or 'sclerotic', due to condensation on the lens nucleus and the deposition of brown pigment within the lens. In its advanced stages, it is called a brunescent cataract. In early stages, an increase in sclerosis may cause an increase in refractive index of the lens. This causes a myopic shift (lenticular shift) that decreases hyperopia and enables presbyopic patients to see at near without reading glasses. This is only temporary and is called second sight.
are due to the lens cortex (outer layer) becoming opaque. They occur when changes in the fluid contained in the periphery of the lens cause fissuring. When these cataracts are viewed through an ophthalmoscope, or other magnification system, the appearance is similar to white spokes of a wheel. Symptoms often include problems with glare and light scatter at night.
are cloudy at the back of the lens adjacent to the capsule (or bag) in which the lens sits. Because light becomes more focused toward the back of the lens, they can cause disproportionate symptoms for their size.
An immature cataract has some transparent protein, but with a mature cataract, all the lens protein is opaque. In a hypermature or Morgagnian cataract, the lens proteins have become liquid. Congenital cataract, which may be detected in adulthood, has a different classification and includes lamellar, polar, and sutural cataracts.
Cataracts can be classified by using the lens opacities classification system LOCS III. In this system, cataracts are classified based on type as nuclear, cortical, or posterior. The cataracts are further classified based on severity on a scale from 1 to 5. The LOCS III system is highly reproducible.
Prevention
Risk factors such as UVB exposure and smoking can be addressed. Although no means of preventing cataracts has been scientifically proven, wearing sunglasses that block ultraviolet light may slow their development. While adequate intake of vitamins A, C, and E may protect against the risk of cataracts, clinical trials have shown no benefit from supplements, although the evidence is mixed, but weakly positive, for a potential protective effect of the carotenoids, lutein and zeaxanthin.
Treatment
Surgical
The appropriateness of surgery depends on a person's particular functional and visual needs and other risk factors. Cataract removal can be performed at any stage and no longer requires ripening of the lens. Surgery is usually outpatient and usually performed using local anesthesia. About 9 of 10 patients can achieve a corrected vision of 20/40 or better after surgery.
Several recent evaluations found that cataract surgery can meet expectations only when significant functional impairment due to cataracts exists before surgery. Visual function estimates such as VF-14 have been found to give more realistic estimates than visual acuity testing alone. In some developed countries, a trend to overuse cataract surgery has been noted, which may lead to disappointing results.
Phacoemulsification is the most widely used cataract surgery in the developed world. This procedure uses ultrasonic energy to emulsify the cataract lens. Phacoemulsification typically comprises six steps:
Anaesthetic – The eye is numbed with either a subtenon injection around the eye (see: retrobulbar block) or topical anesthetic eye drops. The former also provides paralysis of the eye muscles.
Corneal incision – Two cuts are made at the margin of the clear cornea to allow insertion of instruments into the eye.
Capsulorhexis – A needle or small pair of forceps is used to create a circular hole in the capsule in which the lens sits.
Phacoemulsification – A handheld ultrasonic probe is used to break up and emulsify the lens into liquid using the energy of ultrasound waves. The resulting 'emulsion' is sucked away.
Irrigation and aspiration – The cortex, which is the soft outer layer of the cataract, is aspirated or sucked away. Fluid removed is continually replaced with a saline solution to prevent collapse of the structure of the anterior chamber (the front part of the eye).
Lens insertion – A plastic, foldable lens is inserted into the capsular bag that formerly contained the natural lens. Some surgeons also inject an antibiotic into the eye to reduce the risk of infection. The final step is to inject salt water into the corneal wounds to cause the area to swell and seal the incision.
A Cochrane review found little to no difference in visual acuity as a function of the size of incisions made for phacoemulsification in the range from ≤ 1.5 mm to 3.0 mm. Extracapsular cataract extraction (ECCE) consists of removing the lens manually, but leaving the majority of the capsule intact. The lens is expressed through a 10- to 12-mm incision which is closed with sutures at the end of surgery. ECCE is less frequently performed than phacoemulsification, but can be useful when dealing with very hard cataracts or other situations where emulsification is problematic. Manual small incision cataract surgery (MSICS) has evolved from ECCE. In MSICS, the lens is removed through a self-sealing scleral tunnel wound in the sclera which, ideally, is watertight and does not require suturing. Although "small", the incision is still markedly larger than the portal in phacoemulsification. This surgery is increasingly popular in the developing world where access to phacoemulsification is still limited.
Intracapsular cataract extraction (ICCE) is rarely performed. The lens and surrounding capsule are removed in one piece through a large incision while pressure is applied to the vitreous membrane. The surgery has a high rate of complications.
Prognosis
Postoperative care
The postoperative recovery period (after removing the cataract) is usually short. The patient is usually ambulatory on the day of surgery, but is advised to move cautiously and avoid straining or heavy lifting for about a month. The eye is usually patched on the day of surgery and use of an eye shield at night is often suggested for several days after surgery.
In all types of surgery, the cataractous lens is removed and replaced with an artificial lens, known as an intraocular lens, which stays in the eye permanently. Intraocular lenses are usually monofocal, correcting for either distance or near vision. Multifocal lenses may be implanted to improve near and distance vision simultaneously, but these lenses may increase the chance of unsatisfactory vision.
Complications
Serious complications of cataract surgery include retinal detachment and endophthalmitis. In both cases, patients notice a sudden decrease in vision. In endophthalmitis, patients often describe pain. Retinal detachment frequently presents with unilateral visual field defects, blurring of vision, flashes of light, or floating spots.
The risk of retinal detachment was estimated as about 0.4% within 5.5 years, corresponding to a 2.3-fold risk increase compared to naturally expected incidence, with older studies reporting a substantially higher risk. The incidence is increasing over time in a somewhat linear manner, and the risk increase lasts for at least 20 years after the procedure. Particular risk factors are younger age, male sex, longer axial length, and complications during surgery. In the highest risk group of patients, the incidence of pseudophakic retinal detachment may be as high as 20%.
The risk of endophthalmitis occurring after surgery is less than one in 1000.
Corneal edema and cystoid macular edema are less serious but more common, and occur because of persistent swelling at the front of the eye in corneal edema or back of the eye in cystoid macular edema. They are normally the result of excessive inflammation following surgery, and in both cases, patients may notice blurred, foggy vision. They normally improve with time and with application of anti-inflammatory drops. The risk of either occurring is around one in 100. It is unclear whether NSAIDs or corticosteroids are superior at reducing postoperative inflammation.
Posterior capsular opacification, also known as after-cataract, is a condition in which months or years after successful cataract surgery, vision deteriorates or problems with glare and light scattering recur, usually due to thickening of the back or posterior capsule surrounding the implanted lens, so-called 'posterior lens capsule opacification'. Growth of natural lens cells remaining after the natural lens was removed may be the cause, and the younger the patient, the greater the chance of this occurring. Management involves cutting a small, circular area in the posterior capsule with targeted beams of energy from a laser, called Nd:YAG laser capsulotomy, after the type of laser used. The laser can be aimed very accurately, and the small part of the capsule which is cut falls harmlessly to the bottom of the inside of the eye. This procedure leaves sufficient capsule to hold the lens in place, but removes enough to allow light to pass directly through to the retina. Serious side effects are rare. Posterior capsular opacification is common and occurs following up to one in four operations, but these rates are decreasing following the introduction of modern intraocular lenses together with a better understanding of the causes.
Vitreous touch syndrome is a possible complication of intracapsular cataract extraction.
Epidemiology
Age-related cataracts are responsible for 51% of world blindness, about 20 million people. Globally, cataracts cause moderate to severe disability in 53.8 million (2004), 52.2 million of whom are in low and middle income countries.
In many countries, surgical services are inadequate, and cataracts remain the leading cause of blindness. Even where surgical services are available, low vision associated with cataracts may still be prevalent as a result of long waits for, and barriers to, surgery, such as cost, lack of information and transportation problems.
In the United States, age-related lens changes have been reported in 42% between the ages of 52 and 64, 60% between the ages 65 and 74, and 91% between the ages of 75 and 85. Cataracts affect nearly 22 million Americans age 40 and older. By age 80, more than half of all Americans have cataracts. Direct medical costs for cataract treatment are estimated at $6.8 billion annually.
In the eastern Mediterranean region, cataracts are responsible for over 51% of blindness. Access to eye care in many countries in this region is limited. Childhood-related cataracts are responsible for 5–20% of world childhood blindness.
Vision loss due to cataracts increases the risk of dementia in the elderly population, increases the likelihood of falls and road traffic accidents, and by detrimental effects on the quality of life increases mortality.
History
Cataract surgery was first described by the Ayurvedic physician, Suśruta (about 5th century BCE) in Sushruta Samhita in ancient India. Most of the methods mentioned focus on hygiene. Follow-up treatments include bandaging of the eye and covering the eye with warm butter. | Biology and health sciences | Disabilities | Health |
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